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# Using the Magnitude-Squared Coherence for Determining Order-Chaos Transition in a System Governed by Logistic Equation Dynamics

## I Introduction

Nonlinear dynamical methods for analyzing a discrete-time series have been widely used in several fields of scientific research. For characterizing a scalar discrete-time series metric tools of nonlinear dynamical methods, such as dimensions, exponents, and entropies are calculated. Applying these tools to embed a discrete-time series in a phase space to obtain a phase-space portrait of the attractor of the dynamical system is required. In order to make the discrete-time series embedding defining the two embedding parameters which are the embedding dimension, $`D_e`$, and the time delay, $`\tau `$, is required. Several methods have been proposed for determining adequate values of these parameters. When the dimensionality, $`m`$, of the dynamical system is known, it is directly considered that $`D_e=m`$, but it is not the case when we only know a scalar discrete-time series obtained by measuring one of the accessible variables of the state space of a dynamical system. The embedding dimension is usually estimate in accordance with Takens theorem, i.e., $`D_e>2D_2`$, where $`D_2`$ is the correlation dimension of the attractor in the reconstructed phase space and it is required to be also estimate before $`D_e`$ be known. After determining the two embedding parameters, $`D_e`$ and $`\tau `$, one proceeds to calculate estimates of the metric tools. Two important requirements must be considered at this point. One of these is related to the stationarity of the discrete-time series and the other one is related to the length of the data point sequence. Nonlinear dynamical methods applying requires that a discrete-time series to be analyzed be a stationary free-noise sequence during the observation time. For calculating the correlation dimension, $`D_2`$, and the largest Lyapounov exponent, $`\lambda _1`$, Eckmann and Ruelle have suggested a certain minimal value for the length of the experimental data points sequence in order for the results obtained by applying nonlinear dynamical analysis be reliable. It is recognized by several authors that the application of nonlinear dynamical methods implies a high computational cost. In this sense, therefore, any method that allows to obtain any information about the behavior of a dynamical system and that can operate over a relatively short discrete-time series, acquires an important practical value.

For complementing those methods of nonlinear dynamical analysis we shall now try to demonstrate that introducing the concept of magnitude-squared coherence is possible to quantifying the order-chaos transition in a dynamical system described by logistic mapping dynamics. This method is based on the discrete Fourier transform representation of a finite-duration discrete-time series and has the advantage that can be applied to a relatively shorter discrete-time series than that required by nonlinear dynamical methods. The magnitude-squared coherence has been used by several authors for measuring constancy of phase between two or more signals at one frequency. Besides, according to literature the transition from a regular to a chaotic behavior has been studied in many dynamical systems, including the biological type. One of these is a fluid flow, where by calculating nonlinear dynamical parameters the order-chaos transition can be determined.

Van den Bleek and Schouten  determined the order-chaos transition in a fluidized bed by calculating both correlation dimension, $`D_2`$, and Kolmogorov entropy, $`K`$, as a function of the Reynolds number and the superficial gas velocity, indicating two different regimes, order and chaos, and the intermediate region. In a previous work  the DFT representation for qualitatively describing the order-chaos transition in a dynamical system is used. Now we want to make the description of order-chaos transition in a dynamical system governed by the logistic equation dynamics by applying the Welch average periodogram method for calculating the power spectral density and from this, to introduce an inverse average magnitude-squared coherence index for indicating quantitatively order-chaos transition in the dynamical system indicated.

The application of magnitude-squared coherence (MSC) method is reported by several authors. NV Thakor et al presented the power spectral density analysis of electrocardiogram (ECG) waveforms as well as isolated QRS complexes and episodes of noise and artifacts. The power spectral analysis showed that the QRS complex could be separated from other interfering signals. A bandpass filter that maximizes the signal-to-noise ratio would be of use in many monitoring instruments. They calculated the coherence function and, from that, the signal-to-noise ratio. Upon carrying out this analysis on experimentally obtained ECG data, they observed that a bandpass filter with a center frequency of 17 Hz and a Q factor of 5 yields the best signal-to-noise ratio. They estimated the coherence function as given by

$`C_{xy}^2={\displaystyle \frac{|G_{xy}(f).G_{xy}(f)^{}|^2}{G_{xx}(f).G_{yy}(f)}}`$ (1)
where $`G_{xx}(f)`$ is the estimates of the auto-power spectra of the QRS complex, $`G_{yy}(f)`$ is the complete ECG cycle, and $`G_{xy}(f)`$ is the cross power spectrum estimate. For calculating the power spectra they applied the Fourier transform of the windowing signal. To improve the resolution of the FFT algorithm and the power spectra, they filled the original data sets with zeros and formed new 1024 point data sets. S Narayanaswamy et al developed three signal processing tools to identify three postulated mechanisms of isolated or single premature ventricular contraction (PVC) generation. Two hours of continuous ECG recording were digitally obtained from several patients with frequent PVC. They applied the magnitude-squared coherence (MSC) spectrum between the sinus intervals $`x[n]`$ and PVC intervals $`y[n]`$ as given by

$`MSC(f)={\displaystyle \frac{|S_x(f).S_y(f)^{}|}{\sqrt{S_{xx}(f)}.\sqrt{G_{yy}(f)}}}`$ (2)
where $`S_x(f)`$ and $`S_y(f)`$ are the Fourier spectra of $`x[n]`$ and $`y[n]`$ respectively. $`S_y(f)^{}`$ is the complex conjugate of $`S_y(f)`$, and $`S_{xx}(f)`$ and $`S_{yy}(f)`$ are the auto-power spectra of $`x[n]`$ and $`y[n]`$ respectively. The power spectral density of the PVC interval series is defined as

$`P(k)=|X(k)|^2`$ (3)
where $`X(k)`$ is the Fourier transform of $`x[n]`$. In this work, details of the used method are not showed, but they used both the power spectral density of the PVC signal and the MSC estimates of the PVC and sinus intervals.

EG Lovett y KM Ropela showed that magnitude-squared coherence (MSC) of intracardiac electrogram may be used to quantify what may be termed organization. Unfortunately, MSC computations are not meaningful from the surface due to the high correlation among measured body surface potentials. They showed that average magnitude-squared auto-bi-coherence (AMSABC) provides an organizational measure of cardiac rhythms from a single surface lead. They showed that AMSABC also discriminates three ventricular tachyarrthythmia. In this work, rather than measuring consistency of phase between two signals at one frequency (as in MSC), measuring consistency between two frequencies in one signal (as in MSABC) were considered. Haris J. Sih et al applied the MSC function for analyzing the spatial organization of epicardial mapping in the frequency domain. They defined the magnitude-squared coherence (MSC) as a frequency domain measure of the phase consistency between two signals, $`x[n]`$ y $`y[n]`$, as given by

$`MSC_{xy}(f)={\displaystyle \frac{|S_{xy}(f)|^2}{S_{xx}(f).S_{yy}(f)}}`$ (4)
where $`S_{xy}(f)`$ is the cross-power spectrum, and $`S_{xx}(f)`$ and $`S_{yy}(f)`$ are the respective auto-power spectra, MSC is dimensionless and can vary from 0 to 1. A MSC value of 1 at some frequency would indicate a linear relationship or perfect coherence between the two signals at that frequency, while a value of 0 would indicate no relationship or perfect incoherence between the two signals at that frequency. MSC is sensitive to noise and interference, whereas uncorrelated noise in either of the two signals decreases MSC, and correlated interference increases MSC. While the relationship between MSC and correlated interference is difficult to characterize, the influence of noise on MSC can be derived for two real signals, $`u(t)`$ and $`v(t)`$, equal to ideal signals, $`x(t)`$ and $`y(t)`$, plus random, independent, zero-mean noise processes, $`n_1(t)`$ and $`n_2(t)`$. It is easy to show that

$`MSC_{uv}={\displaystyle \frac{MSCxy(f)}{1+\frac{1}{SNR_x(f)}+\frac{1}{SNR_y(f)}+\frac{1}{SNR_x(f)SNR_y(f)}}}`$ (5)
where $`MSC_{uv}`$ is the magnitude-squared coherence between $`u(t)`$ and $`v(t)`$, $`MSC_{xy}`$ is the magnitude-squared coherence between $`x(t)`$ and $`y(t)`$, $`SNR_x(f)`$ is the frequency-dependent signal-to-noise ratio between $`x(t)`$ and $`n_1(t)`$, and $`SNR_y(f)`$ is the frequency-dependent signal-to-noise ratio between $`y(t)`$ and $`n_2(t)`$. As expected, as the signal-to-noise ratios increase, $`MSC_{uv}`$ approaches to $`MSC_{xy}`$. $`MSC`$ can be estimated by replacing $`S_{xy}(f)`$, $`S_{xy}(f)`$, and $`S_{xy}(f)`$ with respective estimates by using the Carter method. This method uses average Fast Fourier Transform of weighted overlapped segments of equal length from the two signals to form the spectral estimates. That is, $`MSC`$ is estimated as given by

$`MSC_{xy}(f)={\displaystyle \frac{|_{i=1}^NX_i(f)Y_i^{}(f)|^2}{|_{i=1}^NX_i(f)X_i^{}(f)_{i=1}^NY_i(f)Y_i^{}(f)|}}`$ (6)
where $`X_i(f)`$ and $`Y_i(f)`$ are the FFT of the $`ith`$ weighted segments from $`x(t)`$ and $`y(t)`$, respectively, $`N`$ is the total number of segments, and $`()`$ denoted the complex conjugate. From this method of estimation, one can interpret $`MSC`$ as a measure of the phase consistency between the two signals. At a single frequency $`f_1`$, the cross-power term in the numerator for one segment can be interpreted as a vector with amplitude $`|X_i(f_1)||Y_i(f_1)|`$ and phase $`\mathrm{}X_i(f_1)\mathrm{}Y_i(f_1)`$. Therefore, the sum of the cross-power terms in the numerator of the $`MSC`$ estimate is equivalent to the vector sum of the cross-power vectors for each segment. If $`\mathrm{}X_i(f_1)\mathrm{}Y_i(f_1)=\mathrm{}X_j(f_1)\mathrm{}Y_j(f_1)`$ for all $`i`$ and $`j`$, that is, if a constant phase relationship from one segment to next exists between $`x(t)`$ and $`y(t)`$ at $`f_1`$, the $`N`$ vector terms are all aligned and the vector sum is at a maximum, as is the $`MSC`$ estimate. If the terms are not constant from segment to segment, the $`N`$ vector terms do not align, the vector sum in the numerator is not at a maximum, and the $`MSC`$ estimate must be less than unity at that frequency. As with all methods of spectrum estimation, the $`MSC`$ estimate has statistical bias and variance, which have been explicitly calculated for the case of Gaussian stationary signals. Another form of bias, known as a bias due to misalignment, has also been shown to exist for this estimator. As an example, for two signals, $`x(t)`$ and $`y(t)=bx(tD)+n(t)`$, where $`n(t)`$ is uncorrelated noise, it has been demonstrated that the estimate of the true $`MSC`$ spectrum will be degraded by a factor of $`(1|D|/T)^2`$ where $`T`$ is the time duration of one windowed segment of data. For example, if $`T`$ were approximately $`0.25s`$, a relative delay of $`0.1s`$ between the two signals could decrease the $`MSC`$ estimate to one-third the true value. It is noted that these results are derived for two linearly related signals, and the effects of this bias due to misalignment on two nonlinearly related signals are yet to be determined. With various types of mapping data, as those of cardiology, of neurology, of geophysics, etc, multiple signals are recorded simultaneously, and often, it is necessary to preserve the physical location of these recordings relative to the underlying activity. Under this conditions, $`n`$ signals, $`x_1(t),x_2(t),,x_n(t)`$, are arranged in a $`1`$, $`2`$, or $`D`$ array, or map. By choosing some meaningful reference signal, $`X_r(t)`$, which can be one of the $`n`$ signals in the map, a corresponding map of $`MSC`$ spectra can be formed. That is, what we have is a $`nMSC`$ spectra set: $`MSC_{x_1x_r}(f)`$, $`MSC_{x_2x_r}(f)`$,$`\mathrm{}`$, $`MSC_{x_nx_r}(f)`$, arranged in the same $`1`$, $`2`$, or $`D`$ map. In the case where the reference is one of the $`n`$ signals in the map, the $`rth`$ $`MSC`$ spectrum, $`MSC_{x_rx_r}(f)`$, is unity for all frequencies. It is noted that the map of $`MSC`$ spectra may indicate how the $`n`$ signals relate to the reference, but they do not imply any relationship between one signal of the $`n`$ signals with any other of the $`n`$ signals. Different maps of $`MSC`$ spectra can be generated using different references. A logical and concise method for displaying these $`MSC`$ spectra is problematic as pointed out by many authors. To simplify the data, one can reduce each $`MSC`$ spectrum to a single number, for example, by averaging $`MSC`$ over frequency.

Finally, we refer to the work done by D Liberati et al, where the joint use of total and partial coherence between pairs of $`EEGs`$ simultaneously recorded in a standard set is shown to enhance what is caused direct correlation between cortical subsystems and what is instead related to the spread of the electromagnetic field. A multi-variable $`AR`$ approach is employed in the computation, giving results even for a very short time window, thus allowing coherence to be investigated at the main cortical latencies of evoked potentials. In particular, when a combined visual and somatosensory stimulation is applied, cortical interactions are captured in the frequency domain. In this work, the total squared coherence $`k`$, as a function of the frequency $`f`$ is defined, for every couple $`j,h`$ of the channels, as the squared cross-spectrum $`P_{jh}`$, normalized to the product of the norm of each auto-spectrum $`P_{jj}`$ and $`P_{hh}`$ as given by

$`k_{jh}(f)={\displaystyle \frac{|P_{jh}(f)|^2}{|P_{jj}(f)|.|P_{hh}(f)|}}`$ (7)
where $`k_{jh}`$ is a real function, ranging from 0 to 1, indicating the amount of correlation between the signals $`j`$ and $`h`$ as a function of the frequency $`f`$. In this work we pretend to demonstrate that it is possible to apply the average Welch periodogram for estimating the power spectral density of a chaotic signal and to calculate the magnitude-squared coherence for determining quantitatively the order-chaos transition in a dynamical system described by the logistic mapping dynamics. The fundamental advantage of this method upon the others is that it can operate on a relatively shorter discrete-time series. By estimating the largest Lyapounov exponent was able to evaluate those results obtained here applying an inverse average magnitude-squared coherence index, defined by averaging $`MSC`$ estimate over the frequency.

## II Methods and Materials

### A Obtaining the discrete-time series set

The logistic equation is an one-dimensional mapping of the real axis in the interval (0,1) to itself, that is a prototype of a nonlinear dynamical system widely used and can be formulated as

$`x[n+1]=rx[n](1x[n]),`$ (8)
and it can also be formulated as

$`x[n+1]=r^n(x[1]{\displaystyle \underset{k=0}{\overset{n}{}}}r^{k+1}x[k+1]);`$ (9)
and with the initial condition $`x[1]`$ specified, one can obtained a time sequence of a given length by a recursive action. The character of the behavior of a data sequence obtained with the iterative process can be modified by conveniently selecting a value of the control parameter $`r`$. In this work a value of the initial condition $`x[1]=0.65`$ was taken, and the interval $`r[2.8,4.0]`$ was selected. A set of evenly spaced values of $`r`$, with a step $`\delta =0.01`$, was taken to produce a family of 120 discrete-time series, each containing $`N=2^{10}`$ samples. The building up of the histogram for each time series allows to have immediately a simple statistical characterization of this time series. In order to make easier further computations data points were organized into a rectangular matrix containing 120 columns, each of which is a time series with the length $`N=1024`$ data points.

### B Determining the largest Lyapounov exponent $`\lambda _1`$

Lyapounov exponents spectrum deals with average exponential divergence or convergence of two neighbor orbits in the phase space. It is common to order the spectrum from largest to shortest exponents: $`\lambda _1,\lambda _2,\lambda _3,\mathrm{}`$. A system containing one or more positive exponents is to be defined as chaotic. For calculating the Lyapounov exponents spectrum some algorithms have been developed, according to the particular situation, being one of the most used algorithms that reported by Wolf et al. From the exponents spectrum, the largest exponent, $`\lambda _1`$, decides the behavior of the dynamical system. For calculating the largest exponent the algorithm proposed by Rosenstein et al can be referred too. Details of these algorithm can be found in referred papers. We determined the largest Lyapounov exponent,$`\lambda _1`$, for each time series in the set, using a professional software. The selection of the time delay, $`\tau `$, and the embedding dimension, $`D_e`$, required for reconstructing the phase space is part of the problem and it is necessary to consider some selecting criteria. This selection can be made by different methods. For selecting $`\tau `$ we may take into account:(a) optimal filling of the phase space, (b) by determining the position of the first local minimum of the data autocorrelation function, (c) by taking the first minimum in the mean mutual information plot, (d) based on the optimal tradeoff between redundance and the irrelevance, and by analyzing the differential equations of the system if they are given. For selecting $`D_e`$ it may be indicated: (a) by applying the Grassberger-Procaccia method, which allows to obtain simultaneously both correlation dimension and embedding dimension, (b) by applying the false nearest neighbors method, and (c) by applying geometric considerations. In this work, a time series obtained from the logistic equation was analyzed by using the method (a) for $`\tau `$ determination and by taking a value $`D_e=3`$, which satisfies the Grassberger-Procaccia algoritm condition $`D_e>2D_2`$, and Takens theory, being $`D_2=0.5`$ a value reported by Hoyer et al for this time series type. Literature refers that for selecting the embedding dimension as

$`D_eD_2+1,`$ (10)
it is sufficient to have a good representation for the attractor in the phase space, reconstructed in delayed coordinates from a scalar time series observed. Any reconstruction made with a value of the embedding dimension less than the minimal value will produce a projection, in that dimension, of the original attractor, and therefore it will be more difficult to understand that projection or it will produce false results when an estimation of a nonlinear dynamical parameter is to be made. It is important to take a value for embedding dimension close to $`2D_2+1`$ in order to avoid the computational cost be excessively high.

### C Estimating the IAMSCI using the Welch average periodogram

It is highly meaningful when dealing with a chaotic system keeping in mind that local sensitivity to a small error is the hallmark of such a system. Such dynamical system as the one described by the logistic equation produces a discrete-time non-chaotic series when the control parameter $`r`$ value satisfies the condition $`r<3.57`$. For such a condition a discrete-time series generated by the recursive process does not almost change if the initial condition is modified and there exists phase consistency between the two discrete-time series at a given frequency. However, when the control parameter value is greater than 3.57, a discrete-time series generated by the logistic equation does meaningfully change if the initial condition is slightly modified and phase consistency between the two discrete-time series at a given frequency, in general, must be lost. This effect may be quantified by applying an inverse average magnitude-squared coherence index. In order to apply the Welch average periodogram method for determining the power density spectrum $`S_x[k]`$ of a discrete-time series one follows the following steps:

(a) Decompose the sequence of N (=1024) data points in M (=9) segments, each one with the same length L (=256), and which may be overlapped P samples. The usual value $`P=0.625L`$ was chosen.

(b) Calculate the DFT representation by FFT algorithm for each segment as given by $`X_m[k]=fft(x_m[n])`$, where $`m=1,2,\mathrm{},M`$, and $`n=0,1,2,\mathrm{},L1`$, and fft represents a function in MatLab for calculating the DFT of a given signal. For each segment a L-FFT was calculated, being L the number of points in the sequence $`X_m[k]`$.

(c)Calculate the periodogram for each segment $`m`$ as given by

$`P_m[k]={\displaystyle \frac{1}{L}}|X_m[k]|^2,m=1,2,\mathrm{},M`$ (11)
(d)Finally, calculate the average periodogram as given by

$`S_x[k]={\displaystyle \frac{1}{M}}{\displaystyle \underset{m=1}{\overset{M}{}}}P_m[k]`$ (12)
where $`k=0,1,2,\mathrm{},L1`$. For a short discrete-time series a rectangular window is recommended. For determining the magnitude-squared coherence sequence the procedure described above is applied to a second signal $`y[n]`$ obtaining the average periodogram $`S_y[k]`$, and the average cross-periodogram $`S_{xy}[k]`$ is also calculated. The magnitude-squared coherence sequence was calculated as given by

$`MSC_{xy}[k]={\displaystyle \frac{S_{xy}[k]S_{xy}^{}[k]}{S_x[k]S_y[k]}}`$ (13)
The $`MSC_{xy}[k]`$ sequence describes the phase consistency between the two signals at a given frequency $`k`$. The mean value of the sequence, $`(MSC_{xy}[k])`$, may be used as an index of global coherence between the two signals. When dealing with a chaotic discrete-time series this index must tend to zero, and when dealing with a deterministic discrete-time series this index must tend toward one. However, it was found that a better index may be defined as given by

$`f=10log_{10}{\displaystyle \frac{1}{MSC_{xy}[k]}}.`$ (14)
It makes easier to compare the results obtained here with those obtained by applying the largest Lyapounov exponent estimate.

## III Results

Figures, from 1 to 4, show a sample of four representative time series and their corresponding magnitude-squared coherence spectra. In order to give a better view of each time series it was only considered 128 data points of each discrete-time series in its plotting, but it was completely considered when its corresponding magnitude-squared coherence was calculated. In either case, the particular values of $`r`$ parameter were indicated. Figure 5 depicts results of plotting the inverse average magnitude-squared coherence $`(1/MSC_{xy}[k])`$, calculated for each time series in the set, as a function of the control parameter $`r`$. And figure 6 shows the results of plotting the inverse average magnitude-squared coherence index, in decibels $`(f=10log_{10}(1/MSC_{xy}[k])`$, calculated for each time series in the set, as a function of the control parameter $`r`$ too. Figure 7 shows results of plotting the largest Lyapounov exponent estimate, $`\lambda _1`$, calculated for each time series in the set, as a function of the control parameter $`r`$. This reference serves as a control for the discussion of the results obtained by applying the inverse average magnitude-squared coherence index method.

## IV Discussion of the results

The application of the DFT representation of a chaotic-type signal gives satisfactory results for determining the order-chaos transition in a system described by the logistic equation dynamics. A critical value for control parameter $`r`$ is reported in literature beyond which a sequence produced by the logistic equation exhibits a chaotic behavior. This threshold value approaches $`r3.57`$. On the other hand, figures from 1 to 4 show the time domain representation of some of the time series and their corresponding magnitude-squared coherence spectrum. Comparing figures 1 ($`r=2.91`$) with 2 ($`r=3.64`$), 3 ($`r=3.75`$), and 4 ($`r=3.84`$), it can be deduced the bifurcation effect and the transition to chaos. Note that in figures 2, 3 and 4 the control parameter satisfies the condition $`r>3.57`$, and it corresponds to a value for a positive largest Lyapounov exponent, as it can be estimated from figure 7. It may be observed in those cases that the magnitude-squared coherence spectrum exhibits values far away from one for several values of the frequency in the $`MSC`$ spectrum. It makes the average magnitude-squared coherence to be less than one, corresponding to a situation for which the discrete-time series becomes incoherent and indicating in this case a chaotic behavior of the dynamical system. This situation can be observed in figures 5 and 6 where either the quantity $`(1/MSC_{xy}[k])`$ as the inverse average $`MSC`$ index, $`(f=10log_{10}(1/MSC_{xy}[k])`$, are greater than $`1`$ and $`0`$, respectively, for $`r>3.57`$. Applying the inverse average magnitude-squared coherence index (IAMSCI) to a discrete-time series in a family derived from an observable in a dynamical system is recommended for quantitative detecting the order-chaos transition, and this can be added to the metric tools of the nonlinear dynamical analysis for complementing. Any method that allows to evaluate the behavior of a dynamical system, and that can be applied to a relatively shorter discrete-time series, must be taken into account as an important practical method.

## V Conclusions

The application of the inverse average magnitude-squared coherence index (IAMSCI), $`f`$, to a discrete-time series obtained from the logistic equation can be done on a chaotic discrete-time series. This discrete-time series has the property that exhibits a high level of loss of coherence, and it can be detected by applying the IAMSCI method. This method does not exert a strong requirement on the length of the experimental data sequence. If one deals with a nonlinear dynamical system which can modify its behavior between order and chaos, a time series produced by that system reflects this change, and the IAMSCI method applying to the discrete-time series allow to determine this change. Any method that permits some quantitative evaluation about the behavior of a dynamical system and that can work over a relatively short time series acquires a practical importance. Besides, this method does not require the phase space reconstruction of the attractor.

## Acknowledgement

We acknowledge Pedro A. Sánchez Fernández from the Department of Foreign Languages for the revision of the English version, and Yoelsy Leiva for the edition. We also thank the Ministry of Higher Education of Cuba for partially financial support of the research.

## VI Figures


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# Tensorial perturbations in an accelerating universe

## 1 Introduction

The inflationary paradigm indicates that we live in a flat universe, but observations from the cosmic microwave background anisotropy detected only $`30\%`$ of the material content necessary for the density parameter to reach the unity. The non-luminous matter around the galaxies is insufficient to complete the deficit of $`70\%`$ of the material content of the universe.

The recent measurements on magnitude and redshift of supernovae type Ia, made independently by Perlmutter et al. (1999) and Riess et al. (1998), indicates that our universe is accelerating. Generally, this acceleration is thought to be due to any kind of repulsive gravity , which can be introduced via a negative pressure in the perfect fluid representation of the universe. So, some candidates to this kind of energy are discarded, namely: neutrinos, kinetic energy, radiation. All described by a positive pressure.

To supply the energetic deficit of the universe perhaps we live in a universe dominate by a cosmological constant  or a scalar field component, denominate quintessence ,, . The quintessence models are an advantageous alternative to the cosmological constant, due the relation with the super symmetric models , the problem of the fine-tuning of the cosmological constant  and supergravity models .

On another hand, a negative pressure can be provided in the perfect fluid energy momentum tensor, taking into account a cosmological particle production. The cosmological scenario with particle production, denominated by open system cosmology (OSC), was introduced by Prigogine et. al  with the intention of solving the problem of entropy content in the universe. Note that, Einstein’s equations can hardly provide an explanation for the origin of the cosmological entropy, since they are purely adiabatic and reversible.

Traditionally the variables describing the cosmological fluid are the energy-density ($`\rho `$) and the themodynamical pressure ($`P_{th}`$). In OSC a supplementary variable, the particle number density, enters into the description via a continuity equation with a source of particles. Naturally, the pressure in the stress energy tensor is reinterpreted and a pressure due to the particle creation is added to the thermodynamical pressure, called pressure creation ($`P_c`$).

Considering the source of particles proportional to the product of the number density and Hubble function, the prediction of the actual value for the rate of the particle production is $`10^{16}nucleons/cm^3/Yr`$ . It is nearly the rate predicted by steady-state model . Otherwise, it is far below the detectable limit, the consequences for the dynamics of the universe can be evaluated. OSC can produce solutions for the very early universe, free of the singularity , . Besides that, it can generate inflationary models without the difficulties associated with supercooling  and lead to a universe sufficiently old as to agree with observations  .

So, we consider the OSC framework an promising background, from the observational point of view, to study cosmological perturbations.

The interest in the study of cosmological tensorial perturbations has increased due to the realization that they contribute to the anisotropy of the cosmic microwave background radiation , , , ,.

Gravitational waves provide a new perspective of the observable universe. The energy density of gravitational waves can affect matter, however weakly, causing density fluctuations, peculiar velocities and tidal distortions, for example. On the another hand, a cosmological background of gravitational waves induces redshift perturbations and angular deviations in light transversing it. The observations, at present, do not distinguish among anisotropies due density perturbations and those produced by gravitational waves at recombination era.

In this work we solve the differential equation for tensorial modes in a synchronous gauge for a cosmological background with particle production.

The redshift fluctuations ($`\stackrel{~}{Z}`$) induced by the gravitational waves in OSC background for an accelerated and non-accelerated universe are calculated.

## 2 Background

### 2.1 Open system cosmology

In this section I attempt to give a summary of OSC basic equations. The universe will be considered homogeneous and isotropic, described by the FRW line element

$$ds^2=R^2(\eta )\{d\eta ^2(dr^2+r^2d\theta ^2+r^2\mathrm{sin}\theta ^2d\varphi ^2)\}.$$
(1)
The Einstein field equations for the space-time described by the line element above and the usual energy momentum tensor of an ideal fluid with an additional pressure  due to the particle production ($`P_c`$) are given by

$`3({\displaystyle \frac{R^{}}{R}})^2=\kappa R^2\rho `$ (2)
$`2{\displaystyle \frac{R^{\prime \prime }}{R}}{\displaystyle \frac{R^{}}{R}}=\kappa R^2P=\kappa R^2(P_{th}+P_c),`$ (3)
where the prime means derivative with respect to the conform time $`\eta `$, and $`\kappa =8\pi G`$, $`\rho `$ is the energy density, $`P_{th}`$ is the thermodynamical pressure, and $`P_c`$ is the creation pressure.

The field equations (2) and (3) are coupled with the balance equation for the particle number density

$$n^{}+n\theta =\mathrm{\Psi },$$
(4)
where $`\theta =3\frac{R^{}}{R}`$ and $`\mathrm{\Psi }`$ is the source of the particle production.

The process of transference of the energy from the gravitational field to the production of particles acts as a source of entropy. Using Gibbs relation, the energy conservation law ($`u^\alpha T_{}^{\alpha \beta }{}_{;\beta }{}^{}=0`$, where $`u^\alpha `$ is the velocity four-vector) and balance equation for the particle number density, it is readily obtained that

$`S_{}^{\alpha }{}_{;\alpha }{}^{}=n\dot{\sigma }+\sigma \mathrm{\Psi }={\displaystyle \frac{P_c\theta }{T}}{\displaystyle \frac{\mu \mathrm{\Psi }}{T}},`$
where $`\sigma `$ is the entropy per particle, $`S^\alpha `$ is the entropy current, $`\mu =\frac{\rho +P_{th}}{n}T\sigma `$ is the chemical potential, $`T`$ is the temperature, and the dot means the usual time derivative. Considering the process as adiabatic ( $`\dot{\sigma }=0`$ ) the pressure creation is given by

$$P_c=\frac{\rho +P_{th}}{n\theta }\mathrm{\Psi }.$$
(5)
Combining equations (2), (3), (4) and the state equation

$`P_{th}=\nu \rho ,`$
where $`\nu `$ is a constant, it follows that

$$RR^{\prime \prime }+[(\frac{3\nu +1}{2})\frac{(\nu +1)\mathrm{\Psi }}{2nH}]R^2=0.$$
(6)
In order to integrate equation (6) we write out an expression for the source $`\mathrm{\Psi }`$. Following Lima et al. , a physically reasonable expression for the particle creation rate is $`\mathrm{\Psi }=3n\beta H`$. To have the characteristic time for matter creation the Hubble time itself, $`\beta `$ must be constant. Beyond this, $`\beta `$ must be a positive constant, otherwise the second law of thermodynamics ($`S_{}^{\alpha }{}_{;\alpha }{}^{}>0`$) is violated. So, integrating (6) the expansion of the universe is governed by the scale factor

$$R=\{R_0(\frac{\eta }{\eta _0})\}^{\frac{2}{3(1\beta )(\nu +1)2}},$$
(7)
where

$$R_0=\frac{3(1\beta )(\nu +1)2}{2}(\frac{N_0^{\nu +1}m}{3})^{\frac{1}{2}}\eta _{0}^{}{}_{}{}^{\frac{2}{3(1\beta )(\nu +1)2}}.$$
(8)
The subscript $`0`$ is relative to the present time, $`m`$ is the mass of produced particle and $`N`$ is the particle number.

The universe that evolves according the scale factor (7) emerges from a singularity. However, in the work of the Prigogine et al. , the source of the particles is proportional to the square of the Hubble function and the universe emerges from an initial instability given by an initial number of particles, instead of a singularity.

### 2.2 Acceleration of the universe

According to data from type Ia supernovae observations, the universe is accelerating. Generally, a negative pressure is responsible for the increase of the expansion velocity of the universe and can be generated from the inclusion of the cosmological term  or using a quintessence component , for example.

An alternative view would be to consider particle creation, which naturally redefine the energy-momentum tensor and could account for this increasing expansion velocity. Taking into account (6), we can write the deacceleration parameter in terms of the particle source function, namely

$$q=\frac{1}{2}\frac{4\pi mG\mathrm{\Psi }}{H^3},$$
(9)
where $`m`$ is the rest mass of the produced particles, and we consider a null value for the curvature. Considering $`\mathrm{\Psi }=3n\beta H`$, we obtain an accelerated universe, expanding with the scale factor (7), if $`\beta >\frac{1}{3}`$ .

## 3 Tensorial modes

The background is described by the line element (7). In order to obtain the perturbed field equations for tensorial modes, we must to return to the original field equations, substituting $`\stackrel{~}{g}_{\mu \nu }=g_{\mu \nu }+h_{\mu \nu }`$, where $`g_{\mu \nu }`$ is the background solution and $`h_{\mu \nu }`$ is a perturbation around $`g_{\mu \nu }`$. The perturbed field equations are

$$\delta R_{\mu \nu }=8\pi G\delta S_{\mu \nu },$$
(10)
where $`S_{\mu \nu }=T_{\mu \nu }\frac{1}{2}g_{\mu \nu }T_{}^{\lambda }{}_{\lambda }{}^{}`$, and $`T_{\mu \nu }`$ is the energy momentum tensor. Therefore, the equation for evolution of the tensorial modes in a synchronous gauge is given by

$$h^{\prime \prime }2\frac{R^{}}{R}h^{}+[n^22(\frac{R^{\prime \prime }}{R}\frac{R^2}{R^2})]h=0,$$
(11)
where the spatial dependence of the perturbed quantities are given in terms of spherical harmonics

$`^2h_{ij}(x,\eta )=h(\eta )n^2Q_{ij}.`$
$`Q_{ij}`$ is tracelles transverse eigenfunction in 3-d spatial section, so that $`^2Q_{ij}=n^2Q_{ij}`$.

Substituting the scale factor (7) in equation (11), we find

$$h_a=C_1\eta ^{k_1}J_w(n\eta ),$$
(12)
$$h_b=C_2\eta ^{k_1}Y_w(n\eta ),$$
(13)
where $`C_1`$ and $`C_2`$ are integrations constants,

$$k_1=\frac{1}{2}\frac{3(\beta 1)(\nu +1)2}{3(\beta 1)(\nu +1)+2}$$
(14)
and

$$w=\frac{3}{2}\frac{\beta (\nu +1)\nu +1}{3\beta (\nu +1)3\nu 1}.$$
(15)
The functions $`J`$ and $`Y`$ are the Bessel functions of the first and second kind respectively.

The evolution of the tensorial modes (12) and (13) will depend of the parameter $`\beta `$. Consequently, will depend on the deacceleration parameter (9).

For an accelerated matter dominated universe ($`\beta >\frac{1}{3},\nu =0`$), we have oscillatory decaying modes (Fig. 1 ).

On the other hand, for a non-accelerated universe we have oscillatory growing modes (Fig. 2 ).

For the radiation dominated era the profiles of the tensorial modes are similar to the matter dominated era.

The integration, for the large scale perturbations ($`n0`$), of the equation (11), furnishes the growing mode:

$$h_1(\eta )=C_1\eta $$
(16)
and a mode dependent of the creation parameter

$$h_2(\eta )=C_2\eta ^{\frac{4}{3\beta (\nu +1)3\nu 1}}.$$
(17)
From observational point of view the large scale perturbation has some advantages. The measure of the anisotropy of cosmic microwave background radiation is well established for $`n0`$, that is, for small value of $`l`$ , where $`l`$ denotes the multipolar order in the expansion of the two point correlation function of temperature

$`C(\theta )=\mathrm{\Sigma }_{l=2}^{\mathrm{}}C_lP_l(\mathrm{cos}(\theta )).`$
Using $`0<\beta <\frac{1}{3}`$ the universe is non-accelerated and the mode (17) will behave as a growing mode.

## 4 Induced redshift

How the propagation of light is affected by the less than perfectly symmetric and smooth universe (here gravitational waves) is our goal in this section. Basically, gravitational waves effects can be put in two categories. Direct ones involve energy density of the gravitational waves affecting matter (for example, density fluctuations, peculiar velocities, tidal distortions) and induced effects that act on the propagation of radiation. We calculate the induction, by a cosmological background of gravitational waves, of redshift perturbations in light transversing it in a universe with particle production. The induced effect on the reshift is given by the relation

$$Z=\frac{R_0}{R}1+\stackrel{~}{Z},$$
(18)
where

$$\stackrel{~}{Z}=\frac{R_0}{2R}_0^\eta \frac{h_{11}}{\eta }𝑑\eta .$$
(19)
Writing $`h_{ij}=F(t)\mathrm{exp}i\stackrel{}{k}\stackrel{}{x}`$, where $`\stackrel{}{x}`$ is the comoving distance and $`\stackrel{}{k}`$ is the comoving wave vector, given by $`k=\frac{2\pi R}{\lambda }`$, and considering the large wavelength perturbations, we obtain the induced redshift for the tensorial perturbations (16) and (17), respectively:

$$\stackrel{~}{Z}_1=C_1\eta ^{\frac{3(1\beta )(1+\nu )4}{3(1\beta )(\nu +1)2}}$$
(20)
and

$$\stackrel{~}{Z}_2=C_2\eta ^{\frac{2}{3(1\beta )(\nu +1)2}}.$$
(21)
The induced redshift is a first order effect in the amplitude perturbations, altering our perception of the universe, not the physical constituents themselves. Note that these perturbations in the redshift are dependents of the creation parameter $`\beta `$, consequently these modes are affected by the acceleration of the universe.

Similarities of this effect with gravitational lensing was point out by E. Linder

## 5 Conclusions

The particle production in the universe results in a reinterpretation of the energy momentum tensor. The ordinary pressure is modified and divided in thermodynamic pressure and a negative pressure due to the creation of particles.

The cosmological particle creation furnishes an alternative model to explain the increasing of the expansion velocity of the universe and furnishes an suitable estimative for the age of the universe. Taking into account the scale factor (7) the Hubble function is given by

$`H={\displaystyle \frac{2}{3(1\beta )(1+\nu )t}}.`$
If we do not violate the dominant energy condition and the second law of thermodynamics the constant $`\beta `$ stay in the interval $`0\beta <1`$. Note that, for this interval of $`\beta `$, the universe will be older than the universe described by standard model without creation ($`\beta =0`$). Consequently, the conflict between the age of the universe and the age of the oldest stars in our galaxy is eliminate in OSC , . So, we consider OSC an promising background from the observational point of view to realize perturbations.

We obtain the tensorial modes, for large wavelength perturbations and modes dependent on the scale of the perturbations, considering the OSC as a background. The behavior of these modes are sensitive to the creation parameter and consequently to the dynamics of the universe. Taking into account $`\beta =0`$ we obtain the usual solution for the ordinary fluid without particle production.

The induced redshift ( $`\stackrel{~}{Z}`$) due to the cosmological gravitational waves are determined for large wavelength perturbations. The induced effects are small, but potentially observable. The induced redshift can be , in principle, used to clarify the dynamics of the expansion of the universe, include for redshift $`z>>1`$.

Analyzing the induced redshift expressions,eq. (20) and eq. (21), we conclude that, if we live in a deaccelerated universe ($`\beta <\frac{1}{3}`$) eq. (20) is a decaying mode and eq. (21) a growing mode. Otherwise, for an accelerated universe ($`\beta >\frac{1}{3}`$) we obtain the inverse behaviour for the induced redshift. So, the induced redshift for more distant objects is greater than the close objects, this seems plausible.

The growing mode for the induced redshift (eq.(21)) can contributes for give us an impression that the observed object has a redshifit greater than the expected in a universe without the background of tensorial perturbations, even than we lived in a deaccelerated universe ($`\beta <\frac{1}{3}`$).

The contribution of gravitational waves, generated in a universe with particle creation, in the anisotropy of the cosmic microwave background radiation is subject for a future study.

## 6 Acknowledgments

I like to acknowledgment to the Brazilian agency (CNPQ) for the financial support.


text:
# Untitled Document

Analyse p-adique et

Représentations Galoisiennes

Jean-Marc Fontaine* Institut Universitaire de France et UMR 8628 du CNRS, Mathématique, Université de Paris-Sud, Bâtiment 425, 91405 ORSAY Cedex, France. E-mail: fontaine@math.u-psud.fr

Abstract

The notion of a $`p`$-adic de Rham representation of the absolute Galois group of a $`p`$-adic field was introduced about twenty years ago (see e.g. \[Fo93\]). Three important results for this theory have been obtained recently: The structure theorem for the almost $`C_p`$-representations, the theorem weakly admissible implies admissible and the theorem de Rham implies potentially semi-stable. The proofs of the first two theorems are closely related to the study of a new kind of analytic groups, the Banach-Colmez spaces and the proof of the third uses deep results on $`p`$-adic differential equations on the Robba ring.

2000 Mathematics Subject Classification: 11F80, 11S25, 12H25, 14G22.

Keywords and Phrases: Galois representations, de Rham representations, Semi-stable representations, $`p`$-adic Banach spaces, $`p`$-adic differential equations.

1. Représentations p-adiques

1.1. – Dans tout ce qui suit, $`K`$ est un corps de caractéristique $`0`$, complet pour une valuation discrète, à corps résiduel $`k`$ parfait de caractéristique $`p>0`$. On choisit une clôture algébrique $`\overline{K}`$ de $`K`$, on note $`C`$ son complété et $`||_p`$ la valeur absolue de $`C`$ normalisée par $`|p|_p=p^1`$. On pose $`G_K=\mathrm{Gal}(\overline{K}/K)`$.

Une représentation banachique (de $`G_K`$) est un espace de Banach $`p`$-adique muni d’une action linéaire et continue de $`G_K`$. Avec comme morphismes les applications $`Q_p`$-linéaires continues $`G_K`$-équivariantes, ces représentations forment une catégorie additive $`Q_p`$-linéaire $`(G_K)`$.

Une $`C`$-représentation (de $`G_K`$) est un $`C`$-espace vectoriel de dimension finie muni d’une action semi-linéaire et continue de $`G_K`$. Lorsque $`k`$ est fini, la catégorie $`\mathrm{Rep}_C(G_K)`$ des $`C`$-représentations s’identifie à une sous-catégorie pleine de $`(G_K)`$ :

PROPOSITION \[Fo00\] . — Supposons $`k`$ fini. Si $`W_1`$ et $`W_2`$ sont des $`C`$-représentations, toute application $`Q_p`$-linéaire continue $`G_K`$-équivariante de $`W_1`$ dans $`W_2`$ est $`C`$-linéaire.

Disons que deux représentations banachiques $`S_1`$ et $`S_2`$ sont presque isomorphes s’il existe un triplet $`(V_1,V_2,\alpha )`$$`V_i`$ est un sous-$`Q_p`$-espace vectoriel de dimension finie de $`S_i`$, stable par $`G_K`$, et où $`\alpha :S_1/V_1S_2/V_2`$ est un isomorphisme (dans $`(G_K)`$). Une presque-$`C`$-représentation (de $`G_K`$) est une représentation banachique qui est presque isomorphe à une $`C`$-représentation. On note $`𝒞(G_K)`$ la sous-catégorie pleine de $`(G_K)`$ dont les objets sont les presque-$`C`$-représentations. Elle contient la catégorie $`\mathrm{Rep}_C(G_K)`$ et la catégorie $`\mathrm{Rep}_{Q_p}(G_K)`$ des représentations $`p`$-adiques de dimension finie (de $`G_K`$) comme sous-catégories pleines.

THÉORÈME A \[Fo02\] . — Supposons $`k`$ fini. La catégorie $`𝒞(G_K)`$ est abélienne. Il existe sur les objets de $`𝒞(G_K)`$ une unique fonction additive $`dh:S(d(S),h(S))N\times Z`$ telle que $`dh(W)=(dim_CW,0)`$ si $`W`$ est une $`C`$-représentation et $`dh(V)=(0,dim_{Q_p}V)`$ si $`V`$ est de dimension finie sur $`Q_p`$.

Si $`S`$ et $`T`$ sont des objets de $`𝒞(G_K)`$, les $`Q_p`$-espaces vectoriels $`\mathrm{Ext}_{𝒞(G_K)}^i(S,T)`$ sont de dimension finie et sont nuls pour $`i\{0,1,2\}`$. On a

$`_{i=0}^2(1)^i\mathrm{dim}_{Q_p}\mathrm{Ext}_{𝒞(G_K)}^i(S,T)=[K:Q_p]h(S)h(T)`$.

1.2. – Soit (1) Voir \[Fo88a\] (resp. \[Fo88b\]) pour plus de détails sur la construction de $`\text{B}_{\text{dR}}`$, $`\text{B}_{\text{cris}}`$ et $`\text{B}_{\text{st}}`$ (resp. sur les représentations p-adiques de de Rham et potentiellement semi-stables). $`FR`$ l’ensemble des suites $`(x^{(n)})_{nN}`$ d’éléments de $`C`$ vérifiant $`(x^{(n+1)})^p=x^{(n)}`$ pour tout $`n`$. Avec les lois

$`(x+y)^{(n)}=lim_m\mathrm{}(x^{(n+m)}+y^{(n+m)})^{p^m}\text{ et }(xy)^{(n)}=x^{(n)}y^{(n)}`$

c’est un corps algébriquement clos de caractéristique $`p`$, complet pour la valeur absolue définie par $`|x|=|x^{(0)}|_p`$ et on note $`R`$ l’anneau de la valuation. Son corps résiduel s’identifie au corps résiduel $`\overline{k}`$ de $`\overline{K}`$. L’anneau $`W(R)`$ des vecteurs de Witt à coefficients dans $`R`$ est intègre. Choisissons $`\epsilon ,\pi R`$ vérifiant $`\epsilon ^{(0)}=1`$, $`\epsilon ^{(1)}1`$ et $`\pi ^{(0)}=p`$ et, pour tout $`aR`$ notons $`[a]=(a,0,0,\mathrm{})`$ son représentant de Teichmüller dans $`W(R)`$. L’application $`\theta :W(R)𝒪_C`$ qui envoie $`(a_0,a_1,\mathrm{})`$ sur $`_{nN}p^na_n^{(n)}`$ est un homomorphisme d’anneaux surjectif dont le noyau est l’idéal principal engendré par $`\xi =[\pi ]p`$. On note encore $`\theta :W(R)[1/p]C`$ l’application déduite en rendant $`p`$ inversible. Rappelons que $`B_{dR}^+=\begin{array}{c}\text{lim}\hfill \\ \hfill \end{array}_{nN}W(R)[1/p]/(\xi ^n)`$ et que le corps $`B_{dR}`$ des périodes $`p`$-adiques est le corps des fractions de $`B_{dR}^+`$. Toute unité $`a`$ de $`R`$ s’écrit d’une manière unique sous la forme $`a=a_0a^+`$ avec $`a_0\overline{k}`$ et $`|a^+1|<1`$, la série $`_{n=1}^+\mathrm{}(1)^{n+1}([a^+]1)^n/n`$ converge dans $`B_{dR}^+`$ vers un élément noté $`\mathrm{log}[a]`$ ; on pose $`t=\mathrm{log}[\epsilon ]`$. On a $`B_{dR}=B_{dR}^+[1/t]`$. On note $`A_{cris}`$ le séparé complété pour la topologie $`p`$-adique de la sous-$`W(R)`$-algèbre de $`W(R)[1/p]`$ engendrée par les $`\xi ^m/m!`$ pour $`mN`$. Alors $`A_{cris}`$ s’identifie à un sous-anneau de $`B_{dR}^+`$ contenant $`t`$ et on pose $`B_{cris}^+=A_{cris}[1/p]`$ et $`B_{cris}=B_{cris}^+[1/t]B_{dR}`$. La série $`_{n=1}^+\mathrm{}(1)^{n+1}\xi ^n/np^n`$ converge dans $`B_{dR}^+`$ vers un élément $`\mathrm{log}[\pi ]=\mathrm{log}([\pi ]/p)`$ et on note $`B_{st}`$ la sous-$`B_{cris}`$-algèbre de $`B_{dR}`$ engendrée par $`\mathrm{log}[\pi ]`$. Pour tout $`bR`$ non nul, il existe $`r,sZ`$, avec $`s1`$ et une unité $`a`$ de $`R`$ tels que $`b^s=\pi ^ra`$ et on pose $`\mathrm{log}[b]=(r\mathrm{log}[\pi ]+\mathrm{log}[a])/s`$. On a $`B_{st}=B_{cris}[\mathrm{log}[b]]`$ dès que $`b`$ n’est pas une unité.

Soit $``$ l’ensemble des extensions finies de $`K`$ contenues dans $`\overline{K}`$. Pour tout $`L`$, on pose $`G_L=\mathrm{Gal}(\overline{K}/L)`$ et on note $`L_0`$ le corps des fractions de l’anneau des vecteurs de Witt à coefficients dans le corps résiduel de $`L`$. Le corps $`\overline{K}`$ se plonge de façon naturelle dans $`B_{dR}^+`$ et l’action de $`G_K`$ s’étend de façon naturelle à $`B_{dR}`$, l’anneau $`B_{st}`$ est stable par $`G_K`$. Pour tout $`L`$, on a $`(B_{dR})^{G_L}=L`$ tandis que $`(B_{st})^{G_L}=L_0`$ et l’application naturelle $`L_{L_0}B_{st}B_{dR}`$ est injective.

Pour toute représentation $`p`$-adique $`V`$ de $`G_K`$ de dimension finie $`h`$ sur $`Q_p`$, on pose $`D_{dR}(V)=(\overline{K}_{Q_p}V)^{G_K}`$, $`D_{st}(V)=(B_{st}_{Q_p}V)^{G_K}`$ et, si $`L`$ est une extension finie de $`K`$ contenue dans $`\overline{K}`$, $`D_{st,L}(V)=(B_{st}_{Q_p}V)^{G_L}`$. On a $`dim_KD_{dR}(V)h`$ et on dit que $`V`$ est de de Rham si on a l’égalité. C’est le cas si $`dim_{K_0}D_{st}(V)=h`$ auquel cas on dit que $`V`$ est semi-stable. C’est aussi le cas s’il existe $`L`$ tel que $`dim_{L_0}D_{st,L}(V)=h`$, auquel cas on dit que $`V`$ est potentiellement semi-stable, ou si l’on veut préciser $`L`$, que $`V`$ est $`L`$-semi-stable.

THÉORÈME B . — Toute représentation $`p`$-adique de $`G_K`$ qui est de de Rham est potentiellement semi-stable. Soit $`\overline{K}B_{st}`$ le plus petit sous-anneau de $`B_{dR}`$ contenant $`\overline{K}`$ et $`B_{st}`$. Ce théorème revient à dire que, pour toute représentation de de Rham $`V`$, l’inclusion $`((\overline{K}B_{st})_{Q_p}V)^{G_K}(B_{dR}_{Q_p}V)^{G_K}`$ est une égalité. Berger \[Be02\] en a ramené la preuve à un résultat sur les équations différentielles $`p`$-adiques, résultat qui a ensuite été prouvé indépendamment par André \[An02\], Kedlaya \[Ke02\] et Mebkhout \[Me02\], voir §3.

L’un des intérêts de ce théorème est que l’on dispose d’une description algébrique explicite de la catégorie des représentations potentiellement semi-stables. Le Frobenius usuel sur $`W(R)`$ s’étend de façon naturelle en un endomorphisme $`\phi `$ de l’anneau $`B_{st}`$ (on a $`\phi t=pt`$ et $`\phi (\mathrm{log}[\pi ])=p\mathrm{log}[\pi ]`$). Il existe une unique $`B_{cris}`$-dérivation $`N:B_{st}B_{st}`$ telle que $`N(\mathrm{log}[\pi ])=1`$. L’action de $`\phi `$ et de $`N`$ commutent à celle de $`G_K`$ et $`N\phi =p\phi N`$.

Soit $`L`$ telle que $`L/K`$ est galoisienne. Pour toute représentation $`p`$-adique $`V`$ de $`G_K`$, $`D_{st,L}(V)`$ est un $`(\phi ,N,\mathrm{Gal}(L/K))`$-module filtré de dimension finie, i.e. c’est un $`L_0`$-espace vectoriel $`D`$ de dimension finie, muni

– de deux applications $`\phi :DD`$, $`N:DD`$, la première semi-linéaire relativement à la restriction de $`\phi `$ à $`L_0`$ et bijective, la deuxième linéaire, vérifiant $`N\phi =p\phi N`$,

– d’une action semi-linéaire de $`\mathrm{Gal}(L/K)`$, commutant à $`\phi `$ et à $`N`$,

– d’une filtration indexée par $`Z`$, décroissante, exhaustive et séparée, du $`K`$-espace vectoriel $`D_K=(L_{L_0}D)^{\mathrm{Gal}(L/K)}`$ (si $`D=D_{st,L}(V)`$, on a $`D_K(B_{dR}_{Q_p}V)^{G_K}`$ et, pour tout $`iZ`$, $`\mathrm{Fil}^iD_K=D_K(B_{dR}^+t^i_{Q_p}V)^{G_K}`$).

On pose $`t_H(D)=_{iZ}i.dim_K\mathrm{Fil}^iD_K/\mathrm{Fil}^{i+1}D_K`$. Si $`D=_{\alpha Q}D_\alpha `$ est la décomposition suivant les pentes du $`\phi `$-isocristal sous-jacent, on pose aussi $`t_N(D)=_{\alpha Q}\alpha .dim_{L_0}D_\alpha `$. On dit que $`D`$ est admissible si

a) on a $`t_H(D)=t_N(D)`$,

b) pour tout sous-$`L_0`$-espace vectoriel $`D^{}`$ de $`D`$, stable par $`\phi ,N`$ et $`\mathrm{Gal}(L/K)`$, on a $`t_H(D^{})t_N(D^{})`$ (on a muni $`D_K^{}D_K`$ de la filtration induite).

THÉORÈME C \[CF00\] . — Soit $`L\overline{K}`$ une extension finie galoisienne de $`K`$.

i) Pour toute représentation $`L`$-semi-stable $`V`$, $`D_{st,L}(V)`$ est admissible.

ii) Le foncteur qui à $`V`$ associe $`D_{st,L}(V)`$ induit une équivalence(2) C’est même une équivalence de catégories tannakiennes, cf. \[Fo88b\]. entre la sous-catégorie pleine $`\mathrm{Rep}_{st,L}(G_K)`$ de $`\mathrm{Rep}_{Q_p}(G_K)`$ dont les objets sont les représentations $`L`$-semi-stables et la catégorie des $`(\phi ,N,\mathrm{Gal}(L/K))`$-modules filtrés admissibles.

Remarque. Il était jusqu’à présent d’usage \[Fo88b\] d’appeler faiblement admissible ce que nous appelons ici admissible. On savait (loc.cit., th.5.6.7) que $`D_{st,L}`$ induit une équivalence entre la catégorie $`\mathrm{Rep}_{st,L}(G_K)`$ et une sous-catégorie pleine de la catégorie des modules filtrés (faiblement) admissibles ; on conjecturait que ce foncteur est essentiellement surjectif et c’est ce qui est prouvé dans \[CF00\].

2. Espaces de Banach-Colmez(3) C’est en cherchant à prouver le théorème C que j’ai été conduit à m’intéresser aux presque C-représentations. C’est Colmez qui a compris que les propriétés dont j’avais besoin provenaient de structures analytiques. Cela nous a permis de prouver le théorème C. Colmez a ensuite étudié plus en détail ces structures analytiques \[Co02\]. Ce que je raconte ici est une interprétation, dans le langage de \[Fo02\], §4, de ces travaux de Colmez et devrait être développé dans \[FP02\].

2.1. – Une $`C`$-algèbre de Banach est une $`C`$-algèbre normée complète $`A`$ ; son spectre maximal est l’ensemble $`\mathrm{Spm}_CA`$ des sections continues $`s:AC`$ du morphisme structural. Si $`fA`$ et $`s\mathrm{Spm}_CA`$, on pose $`f(s)=s(f)`$. Une $`C`$-algèbre spectrale est une $`C`$-algèbre de Banach $`A`$ telle que la norme est la norme spectrale, i.e. telle que, pour tout $`fA`$, $`f=sup_{s\mathrm{Spm}_CA}|f(s)|_p`$ ; dans ce cas, $`\mathrm{Spm}_CA`$ est un espace métrique complet (la distance étant définie par $`d(s_1,s_2)=sup_{f1}|f(s_1)f(s_2)|_p`$). Avec comme morphismes les homomorphismes continus de $`C`$-algèbres, les $`C`$-algèbres spectrales forment une catégorie. La catégorie des variétés spectrales affines sur $`C`$ est la catégorie opposée.

Un groupe spectral commutatif affine sur $`C`$ est un objet en groupes commutatifs dans la catégorie des variétés spectrales affines sur $`C`$. Ces groupes forment, de façon évidente, une catégorie additive qui a des limites projectives finies. Le foncteur qui à un groupe spectral commutatif affine associe le groupe topologique sous-jacent est fidèle. Si $`𝒮=\mathrm{Spm}_CA`$ est un groupe spectral commutatif affine, un sous-groupe spectral du groupe topologique sous-jacent est un sous-groupe $`𝒯`$ qui admet une structure de groupe spectral (nécessairement unique) telle que l’inclusion $`𝒯𝒮`$ est un morphisme de groupes spectraux.

Soit $`S`$ un espace de Banach ($`p`$-adique) et $`𝒮_0`$ la boule unité. Un réseau de $`S`$ est un sous-$`Z_p`$-module $`𝒮`$ qui est tel que l’on peut trouver $`r,sZ`$ vérifiant $`p^r𝒮_0𝒮p^s𝒮_0`$. Il revient au même de dire qu’il existe une norme équivalente à la norme donnée pour laquelle $`𝒮`$ est la boule unité.

Une $`C`$-structure analytique sur $`S`$ est la donnée d’un $`C`$-groupe spectral commutatif affine $`𝒮`$ et d’un homomorphisme continu du groupe topologique sous-jacent à $`𝒮`$ dans $`S`$ dont l’image est un réseau et le noyau un $`Z_p`$-module de type fini. On dit que deux $`C`$-structures analytiques $`𝒮`$ et $`𝒯`$ sur $`S`$ sont équivalentes si $`𝒮\times _S𝒯`$ est un sous-groupe spectral de $`𝒮\times 𝒯`$. Un (espace de) Banach analytique (sur $`C`$) est la donnée d’un espace de Banach muni d’une classe d’équivalence de $`C`$-structures analytiques (on les appelle les structures admissibles de $`S`$). On dit que $`S`$ est effectif s’il existe une structure admissible $`𝒮`$ telle que l’application $`𝒮S`$ est injective.

Un morphisme de Banach analytiques $`f:ST`$ est une application $`Q_p`$-linéaire continue telle qu’il existe des structures admissibles $`𝒮`$ de $`S`$ et $`𝒯`$ de $`T`$ et un morphisme $`𝒮𝒯`$ qui induit $`f`$. Les Banach analytiques forment une catégorie additive $`𝒜_C`$.

Si $`S`$ est un Banach analytique et si $`V`$ est un sous-$`Q_p`$-espace vectoriel de dimension finie, le quotient $`S/V`$ a une structure naturelle de Banach analytique. On dit que deux Banach analytiques $`S_1`$ et $`S_2`$ sont presque isomorphes s’il existe des sous-$`Q_p`$-espaces vectoriels de dimension finie $`V_1`$ de $`S_1`$ et $`V_2`$ de $`S_2`$ et un isomorphisme $`S_1/V_1S_2/V_2`$ (de Banach analytiques).

Le groupe sous-jacent à $`𝒪_C`$ a une structure naturelle de groupe spectral commutatif affine : on a $`𝒪_C=\mathrm{Spm}_CC\{X\}`$$`C\{X\}`$ est l’algèbre de Tate des séries formelles à coefficients dans $`C`$ en l’indéterminée $`X`$ dont le terme général tend vers $`0`$. Ceci fait de $`C`$ un espace de Banach analytique effectif. Un Banach analytique vectoriel est un Banach analytique isomorphe à $`C^d`$ pour un entier $`d`$ convenable. Un espace de Banach-Colmez est un Banach analytique presque vectoriel, i.e. un Banach analytique qui est presque isomorphe à un Banach analytique vectoriel. On note $`𝒞_C`$ la sous-catégorie pleine de $`𝒜_C`$ dont les objets sont les espaces de Banach-Colmez.

PROPOSITION (théorème de Colmez (4) C’est à peu près le résultat principal de \[Co02\]. La définition donnée par Colmez de ce qu’il appelle les Espaces de Banach de dimension finie (avec un E majuscule) est légèrement différente. Il n’est pas très difficile de construire une équivalence entre sa catégorie et la nôtre \[FP02\].) . — La catégorie $`𝒞_C`$ est abélienne et le foncteur d’oubli de $`𝒞_C`$ dans la catégorie des $`Q_p`$-espaces vectoriels est exact et fidèle. Il existe sur les objets de $`𝒞_C`$, une unique fonction additive $`dh:S(d(S),h(S))N\times Z`$ telle que $`dh(C^d)=(d,0)`$ et $`dh(V)=(0,dim_{Q_p}V)`$ si $`V`$ est de dimension finie sur $`Q_p`$.

2.2. – La meilleure façon de comprendre les théorèmes A et C c’est d’utiliser le résultat précédent pour les prouver (5) C’est ainsi que Colmez redémontre le théorème C dans \[Co02\]. Moyennant une preuve un peu plus compliquée, on peut n’utiliser qu’un résultat d’analyticité apparemment moins fort ; c’est ce qu’on fait pour prouver le théorème C dans \[CF00\] et le théorème A dans \[Fo02\].. Lorsque $`k`$ est fini, toute presque-$`C`$-représentation est munie d’unee structure naturelle d’espace de Banach-Colmez ; toute application $`Q_p`$-linéaire continue $`G_K`$-équivariante d’une presque $`C`$-représentation dans une autre est analytique. Le fait que $`𝒞(G_K)`$ est abélienne et l’existence de la fonction $`dh`$ résultent alors du théorème de Colmez.

Le principe de la preuve du théorème C est le suivant : On se ramène facilement au cas semi-stable, i.e. au cas où $`L=K`$. Il s’agit de vérifier que, si $`D`$ est un $`(\phi ,N)`$-module filtré (faiblement) admissible de dimension $`h`$, il existe une représentation $`p`$-adique $`V`$ de dimension $`h`$ telle que $`D_{st,K}(V)`$ soit isomorphe à $`D`$. Une torsion à la Tate permet de supposer que $`\mathrm{Fil}^0D_K=D_K`$. Notons $`V_{st}^{+,0}(D)`$ le $`Q_p`$-espace vectoriel des applications $`K_0`$-linéaires de $`D`$ dans $`B_{st}B_{dR}^+`$ qui commutent à l’action de $`\phi `$ et de $`N`$ et $`V_{st}^{+,1}`$ le quotient du $`K`$-espace vectoriel des applications $`K`$-linéaires de $`D_K`$ dans $`B_{dR}^+`$ par le sous-espace des applications qui sont compatibles avec la filtration. On commence par vérifier que le noyau $`V_{st}^{}(D)`$ de l’application évidente $`\beta :V_{st}^{+,0}(D)V_{st}^{+,1}(D)`$ est un $`Q_p`$-espace vectoriel de dimension finie $`h`$ et que, s’il est de dimension $`h`$, alors la représentation duale $`V_{st}(D)`$ est semi-stable et $`D`$ est isomorphe à $`D_{st}(V_{st}(D))`$. La théorie des espaces de Banach analytiques permet de munir $`V_{st}^{+,0}(D)`$ et $`V_{st}^{+,1}(D)`$ d’une structure d’espace de Banach-Colmez et on a $`dh(V_{st}^{+,0}(D))=(t_N(D),h)`$ tandis que $`dh(V_{st}^{+,1}(D))=(t_H(D),0)`$. Il suffit alors de vérifier que l’application $`\beta `$ est analytique. Comme $`t_H(D)=t_N(D)`$, l’additivité de $`dh`$ implique que $`\beta `$ est surjective et que $`dh(V_{st}^{}(D))=(0,h)`$, ce qui signifie bien que $`dim_{Q_p}V_{st}^{}(D)=h`$.

3. Equations différentielles

3.1. – Soit $`A`$ un anneau commutatif et $`d:A\mathrm{\Omega }_A`$ une dérivation de $`A`$ dans un $`A`$-module $`\mathrm{\Omega }_A`$. Ici, un $`A`$-module à connexion (sous-entendu relativement à $`d`$) est un $`A`$-module libre de rang fini $`𝒟`$ muni d’une application $`:𝒟𝒟\mathrm{\Omega }_A`$ vérifiant la règle de Leibniz. On dit que ce module est trivial s’il est engendré par le sous-groupe $`𝒟_{=0}`$ des sections horizontales.

Pour tout corps $`L`$ de caractéristique $`0`$, complet pour une valuation discrète, notons (cf. par exemple \[Ts98\], §2) $`_{x,L}`$ l’anneau de Robba de $`L`$ (ou ‘’anneau des fonctions analytiques sur une couronne d’épaisseur nulle”), c’est-à-dire l’anneau des séries $`_{nZ}a_nx^n`$ à coefficients dans $`L`$ vérifiant

$`s<1,|a_n|s^n0\text{ si }n+\mathrm{}\text{ et }r<1\text{ tel que }|a_n|r^n0\text{ si }n\mathrm{}`$.

Le sous-anneau $`_{x,L}^{\mathrm{}}`$ de $`_{x,L}`$ des fonctions $`a_nx^n`$ telles que les $`a_n`$ sont bornés est un corps muni d’une valuation discrète (définie par $`|a_nx^n|=\mathrm{sup}|a_n|`$) qui n’est pas complet mais est hensélien. Son corps résiduel s’identifie au corps des séries formelles $`E=k_L((x))`$$`k_L`$ désigne le corps résiduel de $`L`$. Pour toute extension finie séparable $`F`$ de $`E`$, il existe une, unique à isomorphisme unique près, extension non ramifiée $`_F^{\mathrm{}}`$ de $`_{x,L}^{\mathrm{}}`$ de corps résiduel $`F`$ . Posons $`_F=_{x,L}_{_{x,L}^{\mathrm{}}}_F^{\mathrm{}}`$. Si $`k_F`$ désigne le corps résiduel de $`F`$, $`L^{}`$ l’unique extension non ramifiée de $`L`$ de corps résiduel $`k_F`$ et si $`x^{}`$ est un relèvement dans l’anneau des entiers de $`_F^{\mathrm{}}`$ d’une uniformisante de $`F`$, l’anneau $`_F`$ s’identifie à l’anneau de Robba $`_{x^{},L^{}}`$.

Notons $`\mathrm{\Omega }_{_{x,L}}^1`$ le $`_{x,L}`$-module libre de rang $`1`$ de base $`dx`$, solution du problème universel pour les dérivations continues en un sens évident. Les modules à connexion sur l’anneau $`_{x,L}`$ forment une catégorie artinienne. Si $`𝒟`$ est un objet de cette catégorie, on dit qu’il est unipotent si son semi-simplifié est trivial. On dit qu’il est quasi-unipotent s’il existe une extension finie séparable $`F`$ de $`k((X))`$ telle que le module à connexion sur $`_F`$ déduit de $`𝒟`$ par extension des scalaires soit unipotent.

Pour tout $`z`$ dans l’aneeau des entiers de $`_{x,K_0}^{\mathrm{}}`$, il existe un unique endomorphisme continu $`\phi `$ de $`_{x,K_0}`$ qui prolonge le Frobenius absolu sur $`K_0`$ et vérifie $`\phi (x)=x^p+pz`$ ; on appelle Frobenius un tel endomorphisme. Pour un tel $`\phi `$, on note encore $`\phi :\mathrm{\Omega }_{_{x,K_0}}^1\mathrm{\Omega }_{_{x,K_0}}^1`$ l’application induite. Soit $`𝒟`$ un module à connexion sur $`_{x,K_0}`$. Une structure de Frobenius sur $`𝒟`$ consiste en la donnée d’un Frobenius $`\phi `$ sur $`_{x,K_0}`$ et d’une application $`\phi `$-semi-linéaire $`\phi _𝒟:𝒟𝒟`$ commutant à $``$.

THÉORÈME (André, Kedlaya, Mebkhout (6) Crew \[Cr98\] a suggéré que ce théorème pouvait être vrai ; il a été prouvé indépendamment par André \[An02\], Mebkhout \[Me02\] et Kedlaya \[Ke01\]. Pour André comme pour Mebkhout, c’est un cas particulier d’un résultat plus général dont la preuve repose sur la théorie de Christol-Mebkhout \[CM\]. La preuve de Kedlaya est plus directe : elle utilise une classification à la Dieudonné-Manin des modules munis d’un Frobenius pour se ramener à un résultat de Tsuzuki \[Ts98\]. Voir \[Co01\] pour une étude comparative plus détaillée.) . — Tout module à connexion sur $`_{x,K_0}`$ qui admet une structure de Frobenius est quasi-unipotent. Avant de montrer comment Berger \[Be2\] déduit le théorème B de cet énoncé, rappelons quelques résultats de \[Fo00\], \[Fo90\] et \[CC98\] (cf. aussi \[Co98\]). Dans tout ce qui suit, $`V`$ est une représentation $`p`$-adique de $`G_K`$ de dimension finie $`h`$.

3.2. – Soit $`K_{\mathrm{}}`$ le sous-corps de $`\overline{K}`$ engendré sur $`K`$ par les racines de l’unité d’ordre une puissance de $`p`$. Posons $`H_K=\mathrm{Gal}(\overline{K}/K_{\mathrm{}})`$ et $`\mathrm{\Gamma }_K=G_K/H_K`$. En utilisant la théorie de Sen \[Se80\], on montre \[Fo00\] que l’union $`\mathrm{\Delta }_{dR}(V)`$ des sous-$`K_{\mathrm{}}[[t]]`$-modules de type fini de $`(B_{dR}_{Q_p}V)^{H_K}`$ stables par $`\mathrm{\Gamma }_K`$ est un $`K_{\mathrm{}}((t))`$-espace vectoriel de dimension $`h`$ et qu’il existe une unique connexion

$$:\mathrm{\Delta }_{dR}(V)\mathrm{\Delta }_{dR}(V)dt/t$$

qui a la propriété que, pour tout sous-$`K_{\mathrm{}}[[t]]`$-module de type fini $`Y`$ stable par $`G_K`$, tout entier $`r0`$ et tout $`yY`$, il existe un sous-groupe ouvert $`\mathrm{\Gamma }_{r,y}`$ de $`\mathrm{\Gamma }`$ tel que, si $`(y)=_0(y)dt/t`$, alors

$`\gamma (y)\mathrm{exp}(\mathrm{log}\chi (\gamma )._0)(y)(\mathrm{mod}t^rY)`$, pour tout $`\gamma \mathrm{\Gamma }_{r,y}`$.

Cette connexion est régulière : le $`K_{\mathrm{}}[[t]]`$-module $`\mathrm{\Delta }_{dR}^+(V)=(B_{dR}^+V)\mathrm{\Delta }_{dR}(V)`$ est un réseau de $`\mathrm{\Delta }_{dR}(V)`$ vérifiant $`(\mathrm{\Delta }_{dR}^+(V))\mathrm{\Delta }_{dR}^+(V)dt/t`$. On a $`D_{dR}(V)=(\mathrm{\Delta }_{dR}(V))^{\mathrm{\Gamma }_K}`$. L’action de $`\mathrm{\Gamma }_K`$ est discrète sur $`\mathrm{\Delta }_{dR}(V)_{=0}`$ ; on en déduit que $`\mathrm{\Delta }_{dR}(V)_{=0}=K_{\mathrm{}}_KD_{dR}(V)`$ donc que $`V`$ est de de Rham si et seulement si le module à connexion $`\mathrm{\Delta }_{dR}(V)`$ est trivial. Ceci se produit si et seulement s’il existe un réseau (nécessairement unique) $`\mathrm{\Delta }_{dR}^0(V)`$ de $`\mathrm{\Delta }_{dR}(V)`$ vérifiant $`(\mathrm{\Delta }_{dR}^0(V))\mathrm{\Delta }_{dR}^0(V)dt`$.

3.3. – Rappelons brièvement la théorie des $`(\phi ,\mathrm{\Gamma })`$-modules \[Fo80\]. Soit $`𝒪__0`$ l’adhérence dans $`W(FR)`$ de la sous-$`W(k)`$-algèbre engendrée par $`[\epsilon ]`$ et $`1/([\epsilon ]1)`$. C’est un anneau de valuation discrète complet dont l’idéal maximal est engendré par $`p`$ et dont le corps résiduel $`E_0`$ est le corps des séries formelles $`k((\epsilon 1))`$ vu comme sous-corps fermé de $`FR`$. Notons $`𝒪_{\widehat{}^{nr}}`$ le séparé complété pour la topologie $`p`$-adique de l’union de toutes les sous-$`𝒪__0`$-algèbres finies étales de $`𝒪__0`$ contenues dans $`W(FR)`$. C’est un anneau de valuation discrète complet dont le corps résiduel est une clôture séparable $`E^s`$ de $`E_0`$. Son corps des fractions $`\widehat{}^{nr}`$ s’identifie à un sous-corps fermé du corps $`\stackrel{~}{B}=W(FR)[1/p]`$, stable par l’action de $`G_K`$ et par le Frobenius $`\phi `$. Le corps $`_K=(\widehat{}^{nr})^{H_K}`$ est une extension finie non ramifiée du corps des fractions de $`𝒪__0`$. Son corps résiduel $`E_K`$ est une extension finie séparable de $`E_0`$ ; le corps résiduel $`k^{}`$ de $`E_K`$ est celui de $`K_{\mathrm{}}`$.

Alors, $`D(V)=(\widehat{}^{nr}_{Q_p}V)^{H_K}`$ est un $`(\phi ,\mathrm{\Gamma }_K)`$-module sur $`_K`$, i.e. un $`_K`$-espace vectoriel de dimension finie $`D`$ muni d’un Frobenius $`\phi `$-semi-linéaire (que l’on note encore $`\phi `$) et d’une action semi-linéaire continue de $`\mathrm{\Gamma }_K`$ commutant à l’action de $`\phi `$ ; ce $`(\phi ,\mathrm{\Gamma }_K)`$-module est étale, i.e. il existe un $`𝒪__K`$-réseau $`𝒟`$ de $`D`$ tel que $`𝒟`$ est le $`𝒪__K`$-module engendré par $`\phi (𝒟)`$. La correspondance $`VD(V)`$ définit une équivalence entre $`\mathrm{Rep}_{Q_p}(G_K)`$ et la catégorie des $`(\phi ,\mathrm{\Gamma }_K)`$-modules étales. 3.4. – Il n’y a pas de flèche naturelle de $`\widehat{}^{nr}`$ dans $`B_{dR}^+`$, ce qui fait que la comparaison entre $`\mathrm{\Delta }_{dR}(V)`$ et $`D(V)`$ n’est pas si facile. Toutefois, si $`aR`$ est non nul, $`[a]W(R)B_{dR}^+`$ est inversible dans $`B_{dR}^+`$, ce qui permet de voir $`1/[a]=[1/a]`$ comme un élément de $`B_{dR}^+`$. Tout élément de $`\stackrel{~}{B}`$ s’écrit d’une manière et d’une seule sous la forme $`_n\mathrm{}p^n[a_n]`$, avec les $`a_nFR`$ ; notons $`\stackrel{~}{B}_{dR}^+`$ le sous-anneau de $`\stackrel{~}{B}`$ formé des séries de ce type qui convergent dans $`B_{dR}^+`$. L’application $`\stackrel{~}{B}_{dR}^+B_{dR}^+`$ est injective et permet d’identifier $`\stackrel{~}{B}_{dR}^+`$ à une sous-$`W(R)[1/p]`$-algèbre de $`B_{dR}^+`$. Pour tout $`rN`$, posons $`\widehat{}_r^{nr,\mathrm{}}=\widehat{}^{nr}\phi ^r(\stackrel{~}{B}_{dR}^+)`$ et, pour tout $`b\widehat{}_r^{nr,\mathrm{}}`$, notons $`\phi _r(b)`$ l’unique $`c\stackrel{~}{B}_{dR}^+B_{dR}^+`$ tel que $`\phi ^r(c)=b`$. On a $`\widehat{}_r^{nr,\mathrm{}}\widehat{}_{}^{nr,\mathrm{}}{}_{r+1}{}^{}`$ ; soit $`\widehat{}^{nr,\mathrm{}}`$ l’union des $`\widehat{}_r^{nr,\mathrm{}}`$. Alors $`_K^{\mathrm{}}=(\widehat{}^{nr,\mathrm{}})^{H_K}`$ est un sous-corps dense de $`_K`$ stable par $`\phi `$. On pose $`D^{\mathrm{}}(V)=(\widehat{}^{nr,\mathrm{}}_{Q_p}V)^{H_K}`$. On peut le calculer à partir de $`D(V)`$ : c’est l’union des sous-$`_K^{\mathrm{}}`$-espaces vectoriels de dimension finie de $`D(V)`$ stables par $`\phi `$. Le résultat principal de \[CC98\] est que $`V`$ est surconvergente, c’est-à-dire que l’application naturelle $`_KD^{\mathrm{}}(V)D(V)`$ est un isomorphisme.

Pour tout $`rN`$, soit $`_{K,r}^{\mathrm{}}=(\widehat{}_r^{nr,\mathrm{}})^{H_K}`$. Alors $`D_r^{\mathrm{}}(V)=(\widehat{}_r^{nr,\mathrm{}}_{Q_p}V)^{H_K}`$ est aussi le plus grand sous-$`_{K,r}^{\mathrm{}}`$-module $`M`$ de type fini de $`D(V)`$ tel que $`\phi (M)_{K,r+1}^{\mathrm{}}M`$. Pour $`r`$ assez grand, l’application naturelle $`_K^{\mathrm{}}_{_{K,r}^{\mathrm{}}}D_r^{\mathrm{}}(V)D^{\mathrm{}}(V)`$ est un isomorphisme. Lorsqu’il en est ainsi, on a $`\phi _r(_{K,r}^{\mathrm{}})K_{\mathrm{}}[[t]]`$ et

$`\mathrm{\Delta }_{dR}(V)=K_{\mathrm{}}((t))_{_{K,r}^{\mathrm{}}}D_r^{\mathrm{}}(V)\text{ et donc }D_{dR}(V)=(K_{\mathrm{}}((t))_{_{K,r}^{\mathrm{}}}D_r^{\mathrm{}}(V))^{\mathrm{\Gamma }_K}.`$

3.5. – Wach \[Wa96\] a montré comment calculer $`D_{st}(V)`$ à partir de $`D^{\mathrm{}}(V)`$ lorsque $`V`$ est de hauteur finie. C’est Berger \[Be02\] qui a compris comment traiter le cas général : Choisissons un relèvement $`x`$ dans l’anneau des entiers de $`_K^{\mathrm{}}`$ d’une uniformisante de $`E_K`$. Si $`K_0^{}`$ désigne le corps des fractions de $`W(k^{})`$, $`_K^{\mathrm{}}`$ s’identifie précisément au sous-anneau $`_{x,K_0^{}}^{\mathrm{}}`$ de l’anneau de Robba $`_{x,K_0^{}}`$. Ce dernier ne dépend pas du choix de $`x`$ et nous le notons $`_K^{rig}`$ ; il contient $`t=\mathrm{log}[\epsilon ]`$. Si $`L`$ est une extension finie de $`K`$ contenue dans $`\overline{K}`$, le corps $`F=(E^s)^{\mathrm{Gal}(\overline{K}/LK_{\mathrm{}})}`$ est une extension finie de $`E_K`$ et l’anneau $`_L^{rig}`$ s’identifie à l’anneau noté $`_F`$ au §3.1.

Rappelons (§1.2) que si $`u=\mathrm{log}[\epsilon 1]`$, on a $`B_{st}=B_{cris}[u]`$. Les actions de $`\phi `$ et de $`\mathrm{\Gamma }_K`$ s’étendent de façon évidente à $`_K^{rig}`$, à $`_K^{rig}[1/t]`$ et à l’anneau $`_K^{rig}[1/t][u]`$ des polynômes en $`u`$ à coefficients dans $`_K^{rig}[1/t]`$. Berger montre que

$$D_{cris}(V)=(_K^{rig}[1/t]__K^{\mathrm{}}D^{\mathrm{}}(V))^{\mathrm{\Gamma }_K}\text{ et }D_{st}(V)=(_K^{rig}[1/t][u]__K^{\mathrm{}}D^{\mathrm{}}(V))^{\mathrm{\Gamma }_K}$$

(l’action de $`N`$ sur $`D_{st}(V)`$ est la restriction de $`d/duid_{D\mathrm{}(V)}`$).

3.6. – Posons $`D=D_K^{rig}(V)=_K^{rig}[1/t]__K^{\mathrm{}}D^{\mathrm{}}(V)`$. En utilisant l’action de $`\mathrm{\Gamma }_K`$ comme au §3.4, on définit une connexion $`:DDdt/t`$ qui commute à l’action de $`\phi `$. Cette connexion est régulière au sens qu’il existe un sous-$`_K^{rig}`$-module $`D^+`$ de $`D`$, libre de rang $`h`$, stable par $`\phi `$ et vérifiant $`(D^+)D^+dt/t`$ (prendre $`D^+=_K^{rig}__K^{\mathrm{}}D^{\mathrm{}}(V)`$). On vérifie que le $`_K^{rig}`$-module libre $`\mathrm{\Omega }_{_K^{rig}}^1`$ admet $`d[\epsilon ]`$ comme base. Mais $`dt/t=[\epsilon ]^1/td[\epsilon ]`$ et $`t`$ n’est pas inversible dans $`_K^{rig}`$. On déduit alors facilement du théorème d’André-Kedlaya-Mebkhout que $`V`$ est potentiellement semi-stable si et seulement s’il existe un sous-$`_K^{rig}`$-module libre $`D^0`$ de $`D`$, libre de rang $`h`$, stable par $`\phi `$ et vérifiant $`(D^0)D^0dt`$.

Il ne reste plus qu’à construire un tel $`D^0`$ lorsque $`V`$ est de de Rham. Fixons un entier $`r_01`$ suffisamment grand pour que $`D_{r_0}^{\mathrm{}}(V)`$ contienne une base de $`D^{\mathrm{}}(V)`$ sur $`_K^{\mathrm{}}`$ et pour que $`x_{r_0}^{\mathrm{}}`$. Pour tout $`rr_0`$ le sous-anneau $`_{K,r}^{rig}`$ de $`_K^{rig}=_{x,K_0^{}}`$ formé des $`_{nZ}a_nx^n`$ vérifiant

$`s<1,|a_n|s^n0\text{ si }n+\mathrm{}\text{ et }a_n(\epsilon ^r1)^n0\text{ si }n\mathrm{}`$

est stable par $`\mathrm{\Gamma }_K`$ et contient $`_{K,r}^{\mathrm{}}`$. Si $`D_r=D_{K,r}^{rig}(V)=_{K,r}^{rig}_{_{K,r}^{\mathrm{}}}D_r^{\mathrm{}}(V)`$, alors $`D`$ est la réunion croissante des $`D_r`$ et $`\phi (D_r)D_{r+1}`$. L’application $`\phi _r`$ induit un homomorphisme de $`_{K,r}^{rig}`$ dans $`K_{\mathrm{}}((t))`$ et un isomorphisme de $`K_{\mathrm{}}((t))_{_{K,r}^{rig}}D_r`$ sur $`\mathrm{\Delta }_{dR}(V)`$. L’application $`\mathrm{\Phi }_r:D_r\mathrm{\Delta }_{dR}(V)`$ qui envoie $`a`$ sur $`1a`$ est injective.

Soit $`\mathrm{\Delta }_{dR}^0(V)`$ le sous-$`K_{\mathrm{}}((t))`$-module de $`\mathrm{\Delta }_{dR}(V)`$ engendré par les sections horizontales. Pour tout $`rr_0`$, soit $`D_r^0=\{aD_r\mathrm{\Phi }_s(a)\mathrm{\Delta }_{dR}^0(V)\text{ pour tout }sr\}`$. On a $`D_r^0D_{r+1}^0`$ et $`D^0=_{rr_0}D_r^0`$ est un sous-$`_K^{rig}`$-module de $`D`$, stable par $`\phi `$ et vérifiant $`(D^0)D^0dt`$. Si $`V`$ est de de Rham on déduit du fait que $`\mathrm{\Delta }_{dR}^0(V)`$ est un réseau de $`\mathrm{\Delta }_{dR}(V)`$ que $`D^0`$ est libre de rang $`h`$ sur $`_K^{rig}`$. D’où le théorème B.

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