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# THE SPATIAL DISTRIBUTION OF FLUORESCENT H2 EMISSION NEAR T TAU

## 1  Introduction

T Tau is a multiple system of pre-main sequence stars, composed of an optically-visible K0 star, T Tau N, and a heavily-extincted system, T Tau S (Dyck et al., 1982), which is itself binary (Koresko, 2000; Köhler et al., 2000; Duchêne et al., 2002). The T Tau system is surrounded by an infalling envelope of a few thousand AU (Calvet et al., 1994), which is in the process of being disrupted by powerful outflows (Momose et al., 1996). Two bipolar outflows have been identified in the system (Solf, Böhm, & Raga, 1988; Böhm & Solf, 1994). One of the outflows runs NW-SE and has been associated with the embedded pair T Tau S. Analysis of the forbidden line emission indicates that this outflow is poorly collimated (Böhm & Solf, 1994; Solf & Böhm, 1999). Bright arcs of forbidden line emission (Robberto et al., 1995) and near infrared H2 line emission (Herbst et al., 1996, 1997) are found crossing this outflow at scales of $``$ 2” - 14” ($``$ 280 - 2000 AU at the distance of Taurus, 140 pc, Kenyon et al. 1994). Burnham nebula, $``$ 8” south of the system, is associated with the blueshifted lobe of this outflow; it continues in jets detected out to $``$ 0.7 pc (Reipurth et al., 1997). The other outflow runs E-W and is associated with T Tau N. It is less bright in the optical, although high velocity forbidden line emission is detected towards the west, ending in Herbig Haro object HH155 located at the edge of the reflection nebulosity NGC 1555 (Hind’s nebula), $``$ 35” west from the star. 12CO emission associated with the E-W outflow shows blue and redshifted lobes of moderately high velocity (Schuster et al. 1994), surrounded by rings of 13CO emission at lower velocity; these rings have been interpreted as arising at the wall of the cavity in the envelope opened by the E-W outflow (Momose et al. 1996).

The structure of the outflows and molecular gas and dust very near T Tau (inner 2”x2”) is uncertain. Observations in the V, R, and I filters with the WFPC2 camera on the Hubble Space Telescope ($`HST`$) by Stapelfeldt (1998, S98) showed a scattered light structure whose morphology suggests an illuminated outflow cavity in the circumsystem envelope. Spectroscopic mapping in the K (2.01-2.42 $`\mu `$m) and H (1.5-1.8 $`\mu `$m) bands by Kasper et al. (2001) probed the inner 2”$`\times `$2” of the T Tau system; unresolved Br$`\gamma `$ emission at the positions of T Tau N and S was observed, but no infrared H2 emission was detected.

In contrast, space-based observations in the UV of molecular hydrogen have indicated structure on small scales near T Tau. Brown et al. (1981) detected extended H2 fluorescent emission in T Tau N and its adjacent nebula in a low resolution large aperture spectrum from the IUE satellite. Valenti et al. (2000) refined the location of the H2 emission by comparing data from the Goddard High Resolution Spectrograph (GHRS) on $`HST`$ in G140L mode with all the available IUE spectra of T Tau. The H2 line fluxes were found to be larger in the IUE spectra suggesting that the molecular hydrogen emission extended beyond 0.2” from the star, the size of the Small Science Aperture in GHRS. In this paper, we present new data from the Space Telescope Imaging Spectrograph (STIS) on board $`HST`$ which illustrate the spatial distribution of fluorescent H2 emission in the T Tau system, with important implications for the interaction of outflows and circumstellar matter.

## 2 Observations

T Tauri was observed with $`HST`$ on 2000 February 7 in program GO8317, using the STIS MAMA detector with the G140L grating. A 2” long slit with a PA=$`30^{}`$ was employed, and the data were taken in binned pixels mode with an exposure time of 1622 s. The spectral range covered from 1118 to 1715 Å with a dispersion of 0.58 Å per pixel resulting in a resolution of $``$1 Å at 1500 Å, and a plate scale of 0.024” per pixel. Standard CALSTIS pipeline procedures were used to reduce the data.

Figure 1 shows the position-dispersion image. Faint emission is seen extending roughly 1.6” $``$ 240 AU southward from the ultraviolet image of T Tau N. Some extended emission is seen to the north, but to a much smaller extent. The spectra of the extended emission were extracted using standard IRAF/STSDAS/X1D tasks, although the extraction was difficult because the spectra have low signal-to-noise ($``$ 8 for the strongest lines) and because extended structures fill the large aperture in the dispersion direction. For the extraction, the position of the background was taken outside of the range corresponding to the geocoronal Lyman $`\alpha `$ and O I lines and far from the extended emission; the background was the same for all the spectra, averaged over 10 pixels. The signal along the extended emission is weak, roughly 20 times weaker than the stellar spectrum. Line emission accounts for about 50% of the total signal. We attempt to identify the strongest emission features, recognizing that there may be weaker, more diffuse emission that we cannot trace.

Figure 2 shows the spectrum of T Tau N, scaled by a factor of 0.2, and spectra extracted in 10 pixel = 0.24” boxes at various offsets from the star along the slit orientation. Lines from the highly-ionized species Si IV 1393/1402, C IV 1458/1550, and He II 1640 are present in the stellar spectrum; the corresponding emission features in the spectra offset at $`\pm `$ 0.39” are consistent with the expected wings of the instrumental point spread function for the stellar emission. The other line features in the extended nebulosity are consistent with fluorescent H2 emission. As shown in Figure 2, the strongest lines arising from de-excitation of the rotational-vibrational level (v’,J’=1,4) of the electronic state $`{}_{}{}^{1}\mathrm{\Sigma }_{u}^{+}`$ populated by the (1-2)P(5) transition can be identified. There is also a suggestion that the strongest lines of the de-excitation of the (v’,J’=1,7) level, populated by the (1-2)R(6) transition are present.

Table 1 shows the observed fluxes of the fluorescent lines due to the (1-2)P(5) transition in several offset spectra. The fluxes are measured above the estimated continuum. The last column shows the ratios of the line fluxes at $`\lambda 1446`$ and $`\lambda 1504`$ Å, where the fluxes have been corrected for reddening with AV=2.2 using the HD29647 extinction law (Calvet et al., 2003). The ratios are consistent with the optically thin ratio (0.73), within the uncertainties given by their low S/N.

## 3 Spatial distribution

To make the line identifications shown in Figure 2, it is necessary to impose wavelength shifts to the spectra of the extended emission. If these shifts were interpreted as velocities, they would require motions of $``$ 2000 km s-1 or more, much larger than the velocities previously reported in the region ($``$ 45 km s-1; (Böhm & Solf, 1994; Eislöffel & Mundt, 1998). A more likely explanation for the wavelength shifts is that they are due to spatial shifts of the emission regions relative to the star within the large (2”) aperture. In this case, the wavelength shift $`\mathrm{\Delta }\lambda `$ would imply a displacement $`x=s\mathrm{\Delta }\lambda /m`$ from the axis of the slit, where $`s`$ is the plate scale and $`m`$ is the dispersion. All the offsets are measured from the position of T Tau N.

Figure 2 indicates at the right the offsets $`x`$ required to fit the H2 lines at each $`y`$ offset from the star. Two spatially-distinct regions of H2 emission can be identified. One region extending to $`y0.7`$” southwest from the star, which is shifted slightly westwards, and a region covering $`0.7\mathrm{"}<y<1.6\mathrm{"}`$ from the star, which has an eastward displacement. These regions can be discerned in Figure 1. The northern region, labeled by N in the sketch below the spectrum in Figure 1, can be seen as three features extending to a distance of about 0.7” southward from the star, at PA $`210^{}`$, corresponding to the H2 line pairs at 1490/1505 Å, 1431/1446 Å, and 1547/1562 Å, although in the latter case the structure near the star is difficult to see due to the strong stellar C IV 1550 Å emission. Further south, these emission features seem to disappear; however, similar pairs of emission lines can be seen shifted roughly 25 Å blueward, extending nearly vertically in the image to about 1.6” south. This region of emission is labeled S in Figure 1. The N region extends to the north of the star as well, but it is difficult to identify spectral features in it. After taking into account the position angle of the slit, we find that the spatial location of the upper N emission region appears to correspond to the the brightest part of the southern optical scattered light structure imaged by S98, while the lower emission region roughly follows the outer border of the scattered light nebula, although displaced from it.

To obtain more information about the spatial distribution of the emission, we constructed a model and compare the predicted emission with the observed spectrum. If the spatial distribution of the emission at a given offset $`y`$ is given by w(x,y), and the intrinsic emission at each point is $`G_\lambda `$, then the observed flux at wavelength $`\lambda `$ is given by the convolution of w(x,y) and $`G_\lambda `$,

$$F_\lambda (y)=a\underset{i=1}{\overset{N}{}}G_{\lambda \mathrm{\Delta }\lambda _i}w(x(i),y)$$
(1)
where $`a`$ is a scaling constant, x(i) is the spatial shift across the slit from its axis, $`\mathrm{\Delta }\lambda _i`$ is the wavelength shift corresponding to $`x(i)`$, and $`N`$ is the number of pixels across the slit.

We used as template $`G_\lambda `$ the spectrum of HH 43, a low-excitation Herbig-Haro object, which clearly shows the fluorescent lines we have identified (Schwartz, 1983). The spectrum is an average of the three IUE large aperture spectra with the highest exposure times, SWP31828, SWP23881, and SWP24924, taken from the Multimission Archive at the Space Telescope Science Institute, which have been processed under the NEWSIPS extraction pipeline. The combined signal to noise of the resultant HH spectrum is 8-13. Since we are interested in the relative distribution of the H2 intensity, we normalized $`w(x,y)`$ to 1 and scaled $`G_\lambda `$ with the constant a to fit the observed spectrum of the extended emission at offset $`y`$ = -0.39”.

The predicted spectra at several offsets $`y`$ from the star are shown in Figure 3 compared to the observed spectra smoothed to the IUE resolution of 6 Å; the corresponding spatial map of the emission $`w(x,y)`$ is shown in Figure 4.

The H2 emission in the southern part of region N (N1 in Figure 4) appears to follow the brightness and spatial distribution of the scattered light image of S98, with some indication that the H2 is somehow spatially more concentrated than the optical region. The H2 and optical emission are similar from $``$ 0.”45 to $``$ 0.”65 from the star, but around $``$ 0.”45, the H2 emission is $``$ 30% narrower than the optical scattered light. Consistently, the H2 emission in the northern part of region N (N2 in Figure 4) is low, since the north part of the optical emission is two magnitudes fainter than the southern portion at the same distance from T Tau N (S98).

The resulting predicted spectra for region N1 are shown at -0.39” and -0.62” offsets in Figure 3. The spectral agreement is reasonable considering the low S/N of the spectra and the contamination by stellar emission at -0.39”. The principal failing of the model is the lack of predicted emission near 1600-1620 Å, which is identified in the spectrum of T Tau N as a mix of CII an Fe II by Valenti et al. (2000). Fe II emission may arise from the $``$40 $`\mathrm{km}\mathrm{s}^1`$ J-shocks discussed below as sources of Lyman $`\alpha `$ photons for the lower emission region. It is also likely that the broad Lyman $`\alpha `$ emission profile of T Tau produces some of the fluorescent sequences observed in Mira (Wood, Karovska, & Raymond, 2002) but not in low excitation HH objects. Some of these sequences include emission near 1610 Å.

For region S, $`y0.85\mathrm{"}`$, the observed line profiles are too narrow to resolve the emitting regions; we found the best fits by introducing two narrow structures at slightly different positions, labeled S1 and S2 in Figure 4. The agreement between predicted and observed spectra in Figure 3 is reasonably satisfactory except for a feature near 1600-1620 Å in the $`y=0.62\mathrm{"}`$ offset spectrum. We should also emphasize again that it would be difficult for us to detect very spatially-extended H2 emission in our large aperture data.

## 4 Physical conditions

### 4.1 Estimates of Lyman $`\alpha `$ luminosities

The total observed flux in the strongest H2 line we observe, (1-7)P(5), $`\lambda `$= 1504.8 Å, can be written as (cf. Jordan et al. 1978)

$$F_{H_2}=3\times 10^{16}N(2,5)A_{H_2}\overline{J}_\nu \mathrm{erg}\mathrm{cm}^2\mathrm{s}^1,$$
(2)
where $`N`$(2,5) is the column density in the vibrational-rotational level (v”,J”=2,5) of H2, $`A_{H_2}`$ is the H2 emitting area projected in the sky in arcsec2, and $`\overline{J}_\nu `$ is the mean intensity of the Lyman $`\alpha `$ radiation capable of exciting the (1-2)P(5) transition111From now on, the bar notation refers to any quantity $`\overline{f}=\varphi _\nu f_\nu d_\nu `$, with $`\varphi _\nu `$ the line profile of the H2 line, centered at $`\nu _0+\mathrm{\Delta }\nu =\nu _0+\nu _0\mathrm{\Delta }v/c`$ from the Lyman $`\alpha `$ line center at $`\nu _0`$. For the (1-2)P(5) transition, $`\mathrm{\Delta }v=+98\mathrm{k}\mathrm{m}\mathrm{s}^1`$., in $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1\mathrm{Hz}^1`$. We have used the line parameters $`B_{(2,51,4)}=1.976\times 10^8`$ erg-1 cm2 s-1, $`A_{(1,47,5)}=1.969\times 10^8`$ s-1 and $`\mathrm{\Sigma }_jA_{(1,4j)}=1.7125\times 10^9`$ s-1 from Abgrall et al. (2000).

If we assume that the H2 emission region N is irradiated by stellar Lyman $`\alpha `$ and that there is no extinction between the star and the region, then the mean intensity $`\overline{J}_\nu `$ reaching the region is related to the specific intensity $`\overline{I}_\nu `$ leaving the star by

$$\overline{J}_\nu \frac{1}{4\pi }\frac{\overline{I}_\nu A_{Ly\alpha }}{d^2},$$
(3)
where $`A_{Ly\alpha }`$ is the area of the Lyman $`\alpha `$ emitting region on the star and $`d`$ the distance between the star and the site where the H2 molecules are pumped. In turn, $`\overline{I}_\nu `$ can be expressed in terms of the flux $`\overline{F}_{Ly\alpha }^{obs}`$ of Lyman $`\alpha `$ observed at Earth, as $`\overline{I}_\nu =\overline{F}_{Ly\alpha }^{obs}D^2/A_{Ly\alpha }`$ (assuming isotropic emission) where $`D`$ is the distance to the star. Combining these expressions, we can write

$$\overline{J}_\nu 3.3\times 10^9\frac{\overline{F}_{Ly\alpha }^{obs}}{\theta ^2}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1\mathrm{Hz}^1\mathrm{sr}^1,$$
(4)
where $`d/D`$ has been approximated by the apparent angular separation between the star and the H2 emitting region on the sky, $`\theta `$, in units of arcsec in eq.(4).

Inserting in eq. (2) we obtain

$$F_{H_2}10^6\frac{A_{H_2}}{\theta ^2}N(2,5)\overline{F}_{Ly\alpha }^{obs}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1$$
(5)
We estimate a lower limit for $`N`$(2,5) by taking the column density of the outermost layers of the cavity where radiation of Lyman $`\alpha `$ can penetrate, to a depth where $`\tau _\nu 1`$. The actual column density producing the emission may be larger than this if the emitting region is not seen face-on. This may be indeed the case, given the low inclination to the line of sight, $`10^{}`$, of the E-W outflow estimated by Momose et al. (1996). The optical depth for a transition line from an upper level j to a lower level i can be written as $`\tau _\nu =N_iB_{ij}h\nu _{ij}\varphi _\nu (1exp(h\nu _{ij}/kT))/4\pi N_iB_{ij}h\nu _{ij}\varphi _\nu /4\pi `$, for a temperature of a few thousand K, with $`\nu _{ij}`$ the frequency of the transition. We assume that the H2 line is Gaussian in shape, with a Doppler width $`\mathrm{\Delta }\nu _D=\nu _{ij}\sqrt{v_{th}^2+v_{tur}^2}/c`$. The agreement with the spectrum of HH43 (§3) suggests a low excitation situation, so we adopt a temperature T$``$1000 K to determine the thermal velocity $`v_{th}`$, and assume a turbulent velocity $`v_{tur}`$10 km s-1 as expected from an oblique shock interface (§5.1). Assuming the line profile is given by $`\varphi _\nu 1/\pi ^{1/2}\mathrm{\Delta }\nu _D`$, we get a column density $`N(2,5)7\times 10^{14}`$ cm-2 from the condition $`\tau _\nu 1`$ for the (1-2)P(5) transition.

Let $`Q`$ be the ratio between $`\overline{F}_{Ly\alpha }^{obs}`$ and the total observed flux $`F_{Ly\alpha }`$. We cannot use the stellar Lyman $`\alpha `$ line profile in the spectrum because it has been heavily absorbed by circumstellar and interstellar extinction. So, we estimate the value of $`Q`$ assuming that the Lyman $`\alpha `$ profile is similar to that of Mg II 2800 Å, since both are resonant lines with extended wings formed by partial redistribution. We use the high resolution line profile of Mg II k $`\lambda `$ 2796.3 in T Tau N, obtained with STIS/NUV-MAMA and echelle grating E230M in our $`HST`$ program GO8627 (Fig. 5, Saucedo et al., in preparation). At $`\mathrm{\Delta }v`$= 98 km s-1 from the center of the Mg II 2796.3 Å line, we determine $`Q1.4\times 10^{12}`$ Hz-1. As illustration, we show in Figure 5 the H2 line profile $`\varphi _\nu `$ at $`\mathrm{\Delta }v`$= 98 km s-1 used for the calculation of $`\overline{F}_{Ly\alpha }^{obs}`$.

The emitting area $`A_{H_2}`$ can be estimated from the height of the extraction box used to obtain the spectra and the width of the relative intensity map w(x,y) as $``$ 0.2$`\times `$0.24 arcsec2. Taking the distance $`\theta `$ between the star and the H2 region as $``$ 0.5”, an average flux of H2 1504.8 Å $`2.5\times 10^{15}`$ $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ in region N (cf. Table 1), the estimated values of $`N(2,5)`$ and $`Q`$, and eq. (5), we can write for the luminosity of the stellar Lyman $`\alpha `$ line

$$8\times 10^3L_{}L_{Ly\alpha }2.1L_{}$$
(6)
where the limits arise from the uncertainty in the reddening correction appropiate to the H2 lines, since the extinction along the line of sight of the nebulosity may not be the same as the extinction of the star; the lower limit comes from taking the observed H2 flux and the upper limit from dereddening this flux by the extinction to T Tau N, $`A_v=2.2`$ (Calvet et al., 2003). In addition, as discussed above, the actual value of $`N(2,5)`$ is likely to be higher because of projection effects, making the required Lyman $`\alpha `$ luminosity lower than estimates in eq.(6). These limits for the Lyman $`\alpha `$ luminosity are consistent with those obtained by scaling by a factor of 30 the dereddened flux of the CIV 1548 line in our stellar spectra, as suggested by Ardila et al. (2002). With F$`{}_{\mathrm{CIV}}{}^{}9.2\times 10^{12}`$ $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$, the scaled Lyman $`\alpha `$ luminosity would be $``$ 0.35 L.

The estimated maximum Lyman $`\alpha `$ luminosity (eq. ) is $``$ 60 % of the accretion luminosity of T Tau N estimated from the NUV spectrum obtained simultaneously with the FUV spectrum, $`\mathrm{L}_{\mathrm{acc}}=3.2\mathrm{L}_{}`$ (Calvet et al., 2003), so it can be accounted for by accretion energy. We thus conclude that the H2 emission in region N is consistent with stellar Lyman $`\alpha `$ pumping.

While the bright rims seen in region N are clearly illuminated by T Tau N, the southern arcs S1 and S2 appear to be located (in projection) well inside the envelope, roughly tracing the outer edge of the extended emission of S98. Lyman $`\alpha `$ from the visible star will be absorbed by dust in the envelope so it cannot penetrate more than one dust mean free path, $`l34`$ AU, estimated from models of infalling envelopes (Whitney & Hartmann, 1993) with typical envelope parameters from Calvet et al. (1994) ($`\dot{\mathrm{M}}=3\times 10^6\mathrm{M}_{}\mathrm{yr}^1,\mathrm{r}_\mathrm{c}=100\mathrm{AU},\mathrm{M}=2\mathrm{M}_{}`$) and a dust opacity at the Lyman $`\alpha `$ wavelength of 1538 cm2 g-1 (Osorio et al., 2002). It is therefore more likely that the required Lyman $`\alpha `$ emission comes from the shocks themselves, as in the case of low excitation HH objects (Curiel et al., 1995).

If Lyman $`\alpha `$ is locally produced, then $`\overline{J}_\nu =\overline{I}_\nu /2`$, so we can write

$$L_{H_2}=6.8\times 10^6N(2,5)QL_{Ly\alpha }$$
(7)
where $`L_{H_2}=4\pi D^2F`$ is the H2 luminosity. Using similar values of $`N(2,5)`$ and $`Q`$ as in region N, we get that the luminosity in Lyman $`\alpha `$ required to excite the H2 fluorescence in region S is

$$3.4\times 10^4L_{}L_{Ly\alpha }0.1L_{},$$
(8)
where again, the lower limit is given by the observed values of the H2 flux, and the upper limit by assuming that the extinction is the same as towards T Tau N, although there is no argument to support the assumption of an homogeneous extinction in the whole region. These limits are consistent with values of the Lyman $`\alpha `$ luminosity in low excitation HH objects; for instance, Curiel et al. (1995) find $`L_{Ly\alpha }0.13L_{}`$ in HH 47A.

### 4.2 Estimates of temperature and Hydrogen column density

Fluorescence requires a significant population in excited levels which also can lead to infrared H2 emission lines. Thus, the observed upper limits to the infrared emission line fluxes in observations covering region N (Kasper et al., 2001) can provide an additional constraint on molecular column densities.

An upper limit for the 2.12 $`\mu `$m H2 line can be estimated from the equivalent width of 0.15 Å, which is the detection threshold for features in the Kasper et al. (2001) observations. Using the continuum flux of $`3.1\times 10^9\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1\mu \mathrm{m}^1`$, we obtain that the flux at the 2.12 $`\mu `$m H2 line needed for a 2 $`\sigma `$ detection would be $`9.3\times 10^{14}`$ $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$.

The flux at 2.12 $`\mu `$m H2 in the optically thin limit is given by $`F_{2.12}=I_{2.12}A_{H_2}/D^2=N(1,3)A_{31}h\nu _{31}/4\pi A_{H_2}/D^2`$, where $`N(1,3)`$ is the column density of the upper level of the transition. With the Einstein value $`A_{31}=3.66\times 10^7s^1`$ (Gautier et al., 1976) and the upper limit for the flux, we obtain $`N(1,3)2.24\times 10^{18}\mathrm{cm}^2`$.

As mentioned before, we have only a lower limit to the actual N(2,5). With this lower limit and the upper limit for N(1,3), we can obtain a lower limit to the excitation temperature in region N through the Boltzmann relation $`N(1,3)/N(2,5)=1/1.57exp(E_{(1,32,5)}/kT)`$, with $`E_{(1,32,5)}/k6840`$ (Herzberg, 1950). Using our estimates, we obtain $`T`$ 800 K, consistent with a low velocity shock. From Figures 6 and 7 in Jordan et al. (1978), the population of the levels with (v”J”=2,5) in this temperature range is $``$ 12% of the populations with v”=2 and any J”, and in turn, the population with v”=2 is $``$ 0.001% of the total population. So, with $`N(2,5)>`$7$`\times 10^{14}`$ cm-2, the column density of the H2 molecule is N(H2)$`>5.8\times 10^{20}`$ cm-2.

We have no measurements of H2 2.12 $`\mu `$m in region S to help constrain the temperature and total H2 density. Information on regions S1 and S2 can be obtained from the \[O I\], \[N II\] and \[S II\] PV diagrams in Böhm & Solf (1994). From their Figure 5, regions S1 and S2 are located within one of the multiple components identified in their work, component D, which corresponds to the blueshifted lobe of the N-S outflow. The line ratios \[O I\]/\[S II\], \[N II\]/\[O I\] and \[N II\]/\[S II\] around the location of region S are 1.25, 0.35 and 0.28, respectively. Comparison with shock model predictions Hartigan et al. (1987) indicate a shock of velocity v=40 km s-1. For this shock velocity, the ratio of Lyman $`\alpha `$/\[O I\] can be obtained from these models, yielding a value of $``$ 73; so, a rough estimate of the \[O I\] luminosity is $`4.6\times 10^6L_{}L_{[OI]}3.5\times 10^4L_{}`$. Unfortunately, measurement of absolute fluxes for \[O I\] to compare with this prediction are not available (Solf 2002, private communication). Inspection of the relative intensities of \[O I\] along the slit with PA=0 centered on the star in Böhm & Solf (1994) indicates that the region at $``$ 2 arcsec south is $``$ 4-5 times fainter than the star. We can estimate the \[O I\] stellar luminosity from the mass loss rate $`\dot{\mathrm{M}_\mathrm{w}}`$, using the relationship log $`\dot{\mathrm{M}_\mathrm{w}}=4.3+\mathrm{log}(L_{[OI]}/L_{})`$ (Hartigan et al., 1995), assuming that outflow rate scales as $``$ 0.1 the accretion rate (Calvet, 1998), and obtaining the mass accretion rate from the accretion luminosity. This finally leads us to an estimate of $`\mathrm{L}_{[\mathrm{OI}]}2\times 10^4\mathrm{L}_{}`$ for the region, which is consistent with our expected limits of the \[O I\] luminosity.

## 5 Discussion

### 5.1 The northern H2 emitting region

Region N appears to be aligned along the bright rims of the reflection nebulosity around T Tau N, §3; which was plausibly interpreted by S98 as the walls of a cavity driven into the surrounding medium by the outflow from this star. The H2 emission requires some heating to excite molecules into the level which can fluoresce with Lyman $`\alpha `$. At distances of $`50`$ AU where the fluorescent emission is observed, local dust temperatures resulting from heating by the central star(s) are predicted to be $`100`$ K (Calvet et al., 1994), whereas excitation temperatures must be $`800`$ K by the argument of the previous section. The most likely explanation is that shock heating is responsible for the excitation of H2 in these regions, as also suggested by the consistency of the observed spectrum with that of a low-excitation Herbig-Haro object (§3).

The impact of the wind emanating from T Tau N on the molecular medium can naturally explain the needed excitation of H2, and it has already being suggested as a possible explanation for the near infrared H2 bright rims that seem to coincide with the scattered light image at larger scales (Herbst et al., 1997). Böhm & Solf (1994) argued that the E-W outflow is a highly collimated, high velocity $``$ 200 $`\mathrm{kms}^1`$ jet at 0.3” $``$ 40 AU from the star from analysis of the forbidden line emission. Alternatively, as proposed previously by Shu et al. (1995), Shang et al. (1998), and more recently by Shang et al. (2002), the jet results from a density enhancement along the outflow axis in a wide-angle wind. Our observations support the second hypothesis in that a wide-angle wind is required to produce the excitation along the walls of the cavity. Moreover, we expect that an oblique shock would be formed at the interaction of the wide-angle wind and the molecular environment, as envisoned by Cantó (1980). Only the normal component of the velocity would be thermalized, reducing the effective pre-shock velocity. The velocities required to produce a low excitation spectrum as that emitted along the cavity (§3) are $``$ 60 km s-1 (Hartigan et al., 1987); higher velocities would result in high ionization UV lines, such as CIV and Si IV, which are clearly absent in our data. Thus, the observed H2 fluorescent emission at the cavity walls may be the clearest (though indirect) evidence to date for the presence of low-density, wide-opening-angle outflows driving cavities into the molecular medium in star-forming regions.

Observations of other shock-excited features close to T Tau N are needed to help constrain wide-angle wind properties.

### 5.2 The southern H2 emitting regions

The double arcs of region S are located as close as 0”.3 (49 AU) to T Tau S. The shape and orientation of this double arc are very similar to the arcs of H2 2.12 $`\mu m`$ and optical forbidden line emission described by Herbst et al. (1996, 1997), interpreted as terminal shocks of H2, \[S II\] and \[Fe II\] (Herbst et al., 1997; Robberto et al., 1995) in the N-S outflow. Figure 4b shows an enlarged view of the region around the T Tau system, indicating the positions of regions N and S and the shocked structures described by Herbst et al. (1997), as well as the scattered light image of S98. As mentioned in §4, the arcs of region S are located inside the region of forbidden line emission labeled D by Böhm & Solf (1994), very near the centroid. The outlines of this region as well as the position of the centroid are indicated in Figure 4b. The mean velocity of region D is v = -44 km s-1, enough to have a J-shock which could excite the fluorescence as in other low excitation HH objects (Curiel et al., 1995). In the case of HH47A, Curiel et al. (1995) find that 90 % of the UV line emission comes from a postshock distance of $`2\times 10^{14}\mathrm{cm}`$ 13 AU, which we could not resolve. Therefore, the coincidence of the H2 UV emission with the centroid of the \[S II\] emission of component D in Böhm & Solf (1994) is to be expected. More sensitive near-infrared observations to constrain the 2.12 $`\mu `$m emission would provide better constraints on column densities in these regions.

## 6 Summary

We describe the spatial distribution of the H2 fluorescent emission in the inner 2”x2” region around the T Tau system. A northern region of emission coincides with the walls of the cavity of the envelope seen in scattered light. Lyman $`\alpha `$ from T Tau N appears to be strong enough to excite the fluorescence, but local temperatures $``$ 1000 K are required to mantain the fluorescence; shock heating is requred since stellar heating is not sufficient. The required temperatures and the spectrum of the nebulosity, which resembles that of a low excitation HH object, indicate a low velocity shock. The shock may be produced by the interaction of a wide-angle wind emanating from T Tau N and the molecular material around it. This represents the clearest evidence for the presence of wide-angle winds at scales $`<`$ 100 AU from the star. A southern region of emission consists of two arcs, with shape and orientation similar to the arcs of H2 2.12 $`\mu m`$ and forbidden line emission crossing the outflow associated with the T Tau S system. The arcs are near the centroid of the forbidden line emission in the blueshifted lobe of the outflow. The velocity of the outflow around the arcs are consistent with J-shocks, which could power the fluorescence as in low excitation Herbig-Haro objects.

Acknowledgments. We wish to thank K. Stapelfeldt for providing us the optical image of T Tau, M. Kasper for the infrared images and spectra of T Tau N, and K. Böhm and J. Solf for the original long slit spectrograms of T Tau N and its surroundings. We also thank H. Abgrall, for the updated tables of Einstein coefficients and wavelengths for the UV H2 transitions. This work was supported by NASA through grants GO-08206.01-97A and GO-08317.01-97A from the Space Telescope Science Institute, and by NASA Origins of Solar Systems grant NAG5-9670. The non HST ultraviolet data used in this paper were obtained from the Multimission Archive at the Space Telescope Science Institute (MAST). STScI is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5-26555. Support for MAST for non-HST data is provided by the NASA Office of Space Science via grant NAG5-7584, and by other grants and contracts.


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# Spin-orbit induced backflow in neutron matter with auxiliary field diffusion Monte Carlo

## I Introduction

In recent investigations Sarsa et al. (2003); Fantoni et al. (2001) of the ground state and the magnetic properties of neutron matter with modern nuclear interactions of the Urbana-Argonne type Wiringa et al. (1995); Pudliner et al. (1997) good agreement was observed between results obtained with the Auxiliary Field Diffusion Monte Carlo (AFDMC) Schmidt and Fantoni (1999), a calculation by Morales, Jr. et al. (2002), performed with the Variational Chain Summation (VCS) method Akmal and Pandharipande (1997); Akmal et al. (1998), and Brueckner Hartree Fock estimates Baldo et al. (2000); Vidaña et al. (2002).

However, in spite of the overall agreement of the equation of state and the spin susceptibility, marked differences exist between the AFDMC and the VCS calculations, concerning the contribution to the energy due to the spin–orbit component of the two–body interaction. VCS calculations, performed with the Argonne $`v_{18}`$ Wiringa et al. (1995) and the Urbana-IX three–body potential (AU18 Hamiltonian), find large and negative contributions from the cluster terms $`C_{LS}`$ with either spin-orbit correlations in the trial function and/or the spin–orbit potential. For instance, in correspondence with optimal trial functions, $`C_{LS}`$ amounts to -5.8 MeV at a density $`\rho =\rho _0=0.16`$ fm-3 and -12.1 MeV at twice the same density. The AFDMC calculations of Ref. Sarsa et al. (2003), performed with the simplified, but still realistic, version of the Argonne two–body potential, $`v_8^{}`$ Pudliner et al. (1997), plus the Urbana IX three–body interaction ($`AU8^{}`$ Hamiltonian) get energy differences $`\mathrm{\Delta }E_{LS}`$ between the $`AU8^{}`$ Hamiltonians with and without the spin–orbit potential which are small and positive. At $`\rho =\rho _0`$, $`\mathrm{\Delta }E_{LS}=0.2`$ MeV and at $`\rho =2\rho _0`$, $`\mathrm{\Delta }E_{LS}=1.12`$ MeV.

Such a discrepancy cannot be ascribed to differences between the two potentials $`v_{18}`$ and $`v_8^{}`$. It is well known that they provide very close results for the energy per particle Pudliner et al. (1997). In addition, this discrepancy is confirmed by other FHNC/SOC calculations Sarsa et al. (2003), performed by using the approximations of Ref. Wiringa et al. (1988) and the $`AU8^{}`$ Hamiltonian. They give $`\mathrm{\Delta }E_{LS}=3.7`$ MeV at $`\rho _0`$ and $`\mathrm{\Delta }E_{LS}=10.1`$ MeV at $`2\rho _0`$, with $`C_{LS}`$ being $`3.9`$ MeV and $`10.7`$ MeV respectively.

A possible source for this disagreement might be the use of a less than satisfactory guiding function in the AFDMC method 111 In the AFDMC method a path constraint is used to deal with the fermion sign or phase problem. In all the calculations reported here the real part of the guiding function, evaluated at the walker position and spin, is constrained to be positive in analogy with the fixed node approximation for central potentials. For convenience we refer to the effect of the constraint as the nodal structure of the guiding function. . The nodes of the plane wave Slater determinant might be too poor, particularly when the interaction includes a spin–orbit potential, like in $`AU8^{}`$. The results of the FHNC/SOC calculations of Ref. Sarsa et al. (2003) with the $`AU8^{}`$ Hamiltonian and a trial function of the type $`F_6`$ (not containing spin–orbit correlations), give values for $`\mathrm{\Delta }E_{LS}`$ quite close to the AFDMC ones, which may confirm this hypothesis.

To clarify this issue we modify the guiding functions in our quantum Monte Carlo calculations to contain explicit spin–orbit correlations. This is efficiently done by considering orbitals of the spin–backflow form in the Slater determinant, as explained below.

The structure of this paper is as follows. In Sec. II we show the guiding function used. The computational details are given in Sec. III. The results are shown and discussed in Sec. IV. The conclusions and perspectives of the present work can be found in Sec. V.

## II Spin-Orbit Induced Backflow

The $`𝑳𝑺`$ correlation in FHNC/SOC and VCS calculations takes the form

$`F_b(1,2)={\displaystyle \frac{1}{4i}}f_b(r_{12})\left[𝒓_{12}\times \left(\mathbf{}_1\mathbf{}_2\right)\right]\left(𝝈_1+𝝈_2\right)`$ (1)
where $`r_{12}=|𝒓_{12}|`$ and $`𝝈_j`$ are the Pauli matrices for the $`j`$th particle.

By inspection of the cluster terms of $`AU8^{}`$ at the two–body order, one observes that the leading of the spin–orbit terms are those usually denoted as bbc and cbb in FHNC/SOC theory Pandharipande and Wiringa (1979). In the calculations of Ref. Sarsa et al. (2003) these terms are responsible for $`80\%`$ of the total contribution from all the $`C_{LS}^{(2)}`$ terms in the density range $`(3/4)\rho _0\rho (5/2)\rho _0`$. One can easily prove that exactly the same expressions of the cbb and bbc terms are obtained by the following simplified $`𝑳𝑺`$ correlation.

$`\stackrel{~}{F}_b(1,2)={\displaystyle \frac{1}{4i}}f_b(r_{12})\left[𝒓_{12}\left(\mathbf{}_1\times 𝝈_1\mathbf{}_2\times 𝝈_2\right)\right].`$ (2)
It is found that, at the two–body level of the FHNC/SOC theory, $`\stackrel{~}{F}_b`$ leads to an energy which is only $`10\%`$ different from that obtained with $`F_b`$.

The important feature of the $`𝑳𝑺`$ correlation of Eq. (2) for AFDMC calculations is that, similarly to the case of standard backflow Schmidt et al. (1981), it can be implemented in quantum Monte Carlo simulations by substituting the plane wave orbitals of the Slater function with the following spin–backflow ones

$`\mathrm{exp}(i𝒌𝒓_j)`$ (3)
$`\varphi _𝒌(j)=\mathrm{exp}\left(i𝒌𝒓_j+{\displaystyle \frac{\beta }{2}}{\displaystyle \underset{kj}{}}f_b(r_{jk})(𝒓_{jk}\times 𝒌)𝝈_j\right),`$
where $`\beta `$ is a spin–orbit strength parameter. For $`\beta =1`$ there is a direct correspondence of $`\varphi _𝒌(j)`$ with $`\stackrel{~}{F}_b`$. This spin–backflow ansatz can also be used in the GFMC simulations of small nuclear systems Pudliner et al. (1997); Pieper et al. (2002); Carlson et al. (2003), to include $`𝑳𝑺`$ correlations.

We present and discuss, in the following, the results obtained in AFDMC simulations of neutron matter energy with the $`AU8^{}`$ Hamiltonian and the nodal surface of the spin–backflow Slater function. We will show that these nodes serve to lower the AFDMC energy per particle of neutron matter by a sizable amount, which however is too small to solve the CVS and AFDMC discrepancy, particularly at higher densities.

## III Calculation Details

We describe the neutron system by the non relativistic Hamiltonian

$`H`$ $`=`$ $`T+V_2+V_3`$ (4)
$`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m}}{\displaystyle \underset{j=1,N}{}}_j^2+{\displaystyle \underset{j<k}{}}v_{jk}+{\displaystyle \underset{j<k<l}{}}V_{jkl},`$
where $`m`$ is assumed to be the average of the neutron and proton masses and $`\mathrm{}^2/m=41.47108`$ MeV fm2; the two-body and three-body potentials, $`v_{ij}`$ and $`V_{jkl}`$, are the Argonne $`v_8^{}`$ and the the Urbana IX potentials Pudliner et al. (1997). The three-body contributions to the energy are large, particularly at high densities, and cannot be neglected even in a survey calculation like ours.

In neutron matter this interaction can be written in terms of four components:

$`v_{ij}`$ $`=`$ $`{\displaystyle \underset{p=1}{\overset{4}{}}}v_p(r_{ij})𝒪^p(ij)`$ (5)
$`𝒪^p(ij)`$ $`=`$ $`1,𝝈_i𝝈_j,S_{ij},𝑳𝑺,`$ (6)
where $`S_{ij}`$ and $`𝑳𝑺`$ are the usual tensor and spin–orbit operators. The functions $`v_p(r_{ij})`$ can be found on Ref. Pudliner et al. (1997) and also on Ref. Sarsa et al. (2003).

The Urbana-IX three-body interaction is given by the sum

$`V_{jkl}=V_{jkl}^{SI}+V_{jkl}^{SD},`$ (7)
where $`V_{jkl}^{SI}`$ is a spin independent three–body short range part, and the spin–dependent part, $`V_{jkl}^{SD}`$, in neutron matter, reduces to a sum of terms containing only two–body spin operators, with a form and strength that depends on the positions of three particles. Their explicit expressions can be found in Ref. Sarsa et al. (2003).

We have also considered an interaction obtained from $`AU8^{}`$ by dropping the spin–orbit term, and, as in Ref. Sarsa et al. (2003), it is denoted as $`AU6^{}`$.

The AFDMC Schmidt and Fantoni (1999) method used in these calculations has been previously described in detail Sarsa et al. (2003). It allows Monte Carlo simulations to be performed on a relatively large nuclear system at the required accuracy, thanks to the introduction of auxiliary fields, which uncouple the spin–dependent interaction between particles by means of a Hubbard–Stratonovich transformation. While the propagation of the particle coordinates is done as in diffusion Monte Carlo, that of the spin variables, after having sampled the auxiliary fields, results in a rotation of each particle’s spinor.

The guiding function used in this work is given by a Jastrow product correlating an $`N\times N`$ Slater determinant

$`R,S|\mathrm{\Psi }_{JSB}={\displaystyle \underset{j<k}{}}f(r_{jk})\mathrm{\Phi }_B(R,S),`$ (8)
where $`R`$ and $`S`$ denote the set of particle coordinates and of spinors respectively. The space and spin orbitals in $`\mathrm{\Phi }_B(R,S)`$ are given in Eq. (3), and they operate on the spin states in the following way

$`\varphi _𝒌(i)|`$ $`=`$ $`e^{i𝒌𝒓_i}`$ (9)
$`\times `$ $`[\left(\text{cosh}(A_𝒌(i))+\widehat{A}_𝒌^z(i)\text{sinh}(A_𝒌(i))\right)|`$
$`+`$ $`(\widehat{A}_𝒌^x(i)+i\widehat{A}_𝒌^y(i))\text{sinh}(A_𝒌(i))|],`$
where

$`𝑨_𝒌(i)={\displaystyle \frac{\beta }{2}}{\displaystyle \underset{ji}{}}f_b(r_{ij})𝒓_{ij}\times 𝒌`$ (10)
and $`\widehat{A}_𝒌^\alpha (i)`$, $`\{\alpha =x,y,z\}`$ are the components of an unit vector in the direction of $`𝑨_𝒌(i)`$. A similar expression can be obtained for the action over the spin down single particle state. We denote this function by JSB as opposed to JS which stands for a function without explicit backflow correlations, i. e., with simple plane waves in the spatial part of the orbitals in the Slater determinant ($`f_b(r)=0`$). The AFDMC method described in Ref. Sarsa et al. (2003) needs only slight modifications to deal with $`R,S|\mathrm{\Psi }_{JSB}`$ as guiding function. They refer mainly to calculating the kinetic energy part of the local energy, and the gradient of the guiding function in the drift of the walker.

The Jastrow and spin–orbit correlation functions $`f(r)`$ and $`f_b(r)`$ have been taken as the first and fourth components respectively of the FHNC/SOC correlation operator that minimizes the energy per particle of neutron matter at the desired density Wiringa et al. (1988).

## IV Results

We have made our calculations within the full simulation box, and we have taken into account the 26 neighboring boxes in the tabulation of the correlations $`f(r)`$ and $`f_b(r)`$, and of the various components $`v_p(r)`$ of the two–body potential, as described in Ref. Sarsa et al. (2003). Our results are therefore already tail corrected for a Hamiltonian with two–body force only. Tail corrections for the three-body potential were not included. However previous analyses Sarsa et al. (2003) have shown that they are small for systems with 66 neutrons.

Fig. 1 shows the AFDMC energies per particle of 14 neutrons in a periodic simulation cell at $`\rho =2\rho _0`$ for the $`AU6^{}`$ and $`AU8^{}`$ interactions, as a function of the strength parameter $`\beta `$. One can see that, for the $`AU8^{}`$ interaction the energy minimum is around $`\beta =1`$, consistent with the conjecture that $`f_b(r)`$ moves the nodes in the optimal way when the spin–orbit potential is included in the Hamiltonian. We have obtained a lowering of the energy of 1.6 MeV with respect to the case of the JS nodal surface, which corresponds to $`\beta =0`$ in the figure. If we switch off the spin–orbit potential by considering the $`AU6^{}`$ interaction, the minimum is found at $`\beta =0`$, confirming that the spin–backflow nodes are energetically advantageous only in the presence of the spin–orbit component of $`v_8^{}`$.

AFDMC simulations for the $`AU8^{}`$ Hamiltonian with the JSB guiding function have been also carried out for 66 neutrons. As in the case of the 14 neutron system, we have found a minimum of the energy at $`\beta 1`$. The result obtained is compared in Table 1 with those at $`\beta =0`$ for both the $`AU8^{}`$ and the $`AU6^{}`$ interactions. In spite of sizable differences between the energies per particle of the 14 neutron and the 66 neutron systems, the gain in energy $`(E(JSB)E(JS))/N`$ is roughly independent on the number of particles in the box. The large differences between E(14) and E(66) are mainly due to the effect of three–body interaction. It has been shown Sarsa et al. (2003) that the finite size effects on E(66) are rather small.

The dependence of the JSB energy on the density is reported in Table 2 for the 14 neutron system. There is a very weak dependence of $`\mathrm{\Delta }E_{LS}`$ on the density, in contrast with VCS results.

In order to make comparison with recent quantum Monte Carlo calculations performed for 14 neutrons interacting via the $`v_8^{}`$ two–body potential by Carlson et al. (2003), we report in Table 3 the corresponding AFDMC results. The table displays the energy difference $`\mathrm{\Delta }E_{LS}`$ between the energy obtained with $`v_8^{}`$ and $`v_6^{}`$ two–body potential.

The calculations of Ref. Carlson et al. (2003), to which the results reported in Table 3 refer, have been performed with the $`v_8^{}`$ potential cut off at the edge of the box. Therefore, the values of $`\mathrm{\Delta }E_{LS}`$ extracted from there might not be completely comparable with ours, since ours are already tail corrected and have been obtained without introducing any discontinuity in the potential. AFDMC seems to agree reasonably well with GFMC in the constrained path approximation (GFMC–CP). The VCS estimate seems to be too large.

## V Conclusions

In this paper we have proposed a new kind of space–spin orbitals with a spin–backflow form, which is particularly useful for taking into account the spin–orbit interaction of nuclear systems in quantum Monte Carlo simulations. An efficient parameterization of the spin–backflow is obtained from the spin–orbit correlation of FHNC/SOC theory. These spin-backflow orbitals can be conveniently used also in other quantum Monte Carlo calculations, for instance the Green’s Function Monte Carlo simulations of small nucleon systems. The nodal surface provided by this new guiding function is able to decrease the energy about 5 per cent. This amount is not sufficient to solve the spin–orbit discrepancy between the variational chain summation and the AFDMC results, however the AFDMC results are in good agreement with constrained GFMC simulations.

###### Acknowledgements.

SAV wants to thanks the kind hospitality of SF and SISSA where most of this work has been accomplished. Portions of this work were supported by the Italian MIUR-National Research Project 2001025498. AS acknowledges the Spanish Ministerio de Ciencia y Tecnologia for partial support under contract BMF2002-00200


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# Möbius function of coordinate hyperplanes in complex ellipsoids

## 1. Introduction and main results

Let $`G^n`$ be a domain and let $`AG`$. Define the generalized Green function

$$\begin{array}{c}g_G(A,z):=sup\{u(z):u:G[0,1),\mathrm{log}u𝒫𝒮(G),\hfill \\ \hfill _{aA}_{C=C(u,a)>0}_{wG}:u(w)Cwa\},zG,\end{array}$$

and the generalized Möbius function

$$m_G(A,z):=sup\{|f(z)|:f𝒪(G,𝔻),f|_A0\},zG,$$

where $`𝔻`$ denotes the unit disc (cf. \[Jar-Jar-Pfl 2003\]).

The generalized Möbius and Green functions were recently studied by many authors (e.g. \[Lár-Sig 1998\]). Despite various properties proven, effective formulas for $`m_G(A,)`$ and $`g_G(A,)`$ are known only in a few special cases. An effective formula for the generalized Green function in the case of Euclidean ball $`G=𝔹_2^2`$ and union of coordinate lines $`A=\{(z,w)𝔹_2:zw=0\}`$ was given in \[Ngu 2003\].

Let $`n2`$. Let $`p_1,\mathrm{},p_n>0`$. Define the complex ellipsoid with exponents $`p_1,\mathrm{},p_n`$:

$$𝔼=𝔼_{p_1,\mathrm{},p_n}:=\{z^n:\underset{j=1}{\overset{n}{}}|z_j|^{2p_j}<1\}.$$

Observe that $`𝔼_{p_1,\mathrm{},p_n}`$ is convex if and only if $`p_j1/2`$, $`j=1,\mathrm{},n`$ (cf. \[Jar-Pfl 1993\], § 8.4).

For $`G^n`$, $`k\{1,\mathrm{},n\}`$, we consider the sets

$$A=A_{G,k}:=\{zG:z_1\mathrm{}z_k=0\}.$$

We give an effective formula for $`g_𝔼(A_{𝔼,k},)`$ for any complex ellipsoid $`𝔼^n`$ and any $`k\{1,\mathrm{},n\}`$ (Theorem 1.2(a)). This is a generalization of a result in \[Ngu 2003\] (with a much simpler proof). We prove that the $`m_𝔼(A_{𝔼,k},)g_𝔼(A_{𝔼,k},)`$ for the cases where $`𝔼`$ is convex (Theorem 1.2(b)) or where $`k=1`$, $`n=2`$, $`p_21/2`$ (Theorem 1.2(c)). This is not true in general nonconvex case; we prove that $`m_𝔼(A_{𝔼,k},)g_𝔼(A_{𝔼,k},)`$ if there exists a $`j\{k+1,\mathrm{},n\}`$ such that $`p_j<1/2`$ (Theorem 1.2(d)). Partial results for the case $`k=n=2`$ are given (Theorem 1.2(e)).

###### Definition 1.1.

Let $`p_1,\mathrm{},p_n>0`$, $`1kn`$. Put $`𝔼:=𝔼_{p_1,\mathrm{},p_n}`$, $`A:=A_{𝔼,k}`$. Let $`z𝔼`$ be such that the sequence $`(p_j|z_j|^{2p_j})_{j=1}^k`$ is monotonically increasing. For $`s\{1,\mathrm{},k\}`$ define

$$q_s:=\underset{j=1}{\overset{s}{}}(2p_j)^1,r_s(z):=1\underset{j=s+1}{\overset{n}{}}|z_j|^{2p_j},c_s(z):=r_s(z)/q_s.$$

Let $`d:=\mathrm{max}\{s\{1,\mathrm{},k\}:2p_s|z_s|^{2p_s}c_s(z)\}`$. Define

($``$)
$$R_𝔼(A,z)=\underset{j=1}{\overset{d}{}}|z_j|\left(\frac{2p_j}{c_d(z)}\right)^{\frac{1}{2p_j}}.$$

The main results of the paper are the following:

###### Theorem 1.2.

Under the above assumptions we have:

1. $`g_𝔼(A,)R_𝔼(A,)`$,
2. $`m_𝔼(A,)g_𝔼(A,)R_𝔼(A,)`$, for $`p_j1/2`$, $`j=1,\mathrm{},n`$,
3. $`m_𝔼(A,)g_𝔼(A,)R_𝔼(A,)`$, for $`k=1`$, $`n=2`$, $`p_21/2`$,
4. $`m_𝔼(A,)g_𝔼(A,)`$, if there exists a $`j\{k+1,\mathrm{},n\}`$ with $`p_j<1/2`$.
5. $`m_𝔼(A,)g_𝔼(A,)R_𝔼(A,)`$, for $`k=n=2`$, $`p_1p_2`$, and either $`p_21/2`$ or $`8p_1+4p_2(1p_2)>1`$.

Observe that one can obtain the monotonicity of $`(p_j|z_j|^{2p_j})_{j=1}^k`$ for arbitrary $`z𝔼`$ by permutation of coordinates. The final result is a subdivision of $`𝔼`$ into $`2^k1`$ subsets, each with $`R_𝔼(A,)`$ given by a formula of type ($``$).

###### Remark 1.3.

In the case where $`p_1=\mathrm{}=p_n=1`$, the domain $`𝔼`$ is the Euclidean ball $`𝔹_n`$. Theorem 1.2 may then be formulated in the following simpler way.

###### Corollary 1.4.

Assume that $`|z_1|\mathrm{}|z_k|`$ and let

$$d:=\mathrm{max}\left\{s\{1,\mathrm{},k\}:s|z_s|^2+\underset{j=s+1}{\overset{n}{}}|z_j|^21\right\}.$$

Then

$$m_{𝔹_n}(A,z)=g_{𝔹_n}(A,z)=\left(\frac{d}{1_{j=d+1}^n|z_j|^2}\right)^{\frac{d}{2}}\underset{j=1}{\overset{d}{}}|z_j|.$$

###### Remark 1.5.

It is unknown to the author whether $`m_𝔼(A,)g_𝔼(A,)`$ for any $`p_j>0`$, $`j=1,\mathrm{},k`$, $`p_j1/2`$, $`j=k+1,\mathrm{},n`$. However, Theorem 1.2(e) and author’s research (cf. Remark 4.4) indicate that this is true at least in the case $`k=n=2`$.

###### Remark 1.6.

Take $`p_{j,\mathrm{}}+\mathrm{}`$, $`j=1,\mathrm{},n`$. Then

$$G_{\mathrm{}}:=𝔼_{p_{1,\mathrm{}},\mathrm{},p_{n,\mathrm{}}}𝔻^n.$$

Consequently,

$$m_G_{\mathrm{}}(A_{G_{\mathrm{}},k},)=g_G_{\mathrm{}}(A_{G_{\mathrm{}},k},)m_{𝔻^n}(A_{𝔻^n,k},)=g_{𝔻^n}(A_{𝔻^n,k},)$$

(cf. \[Jar-Jar-Pfl 2003\], Property 2.7) and hence

$$m_{𝔻^n}(A_{𝔻^n,k},z)=|z_1|\mathrm{}|z_k|,z𝔻^n.$$

However, we need the formula for $`g_{𝔻^n}(A_{𝔻^n,k},)`$ before we prove Theorem 1.2 (cf. Lemma 2.1).

## 2. Proof of Theorem 1.2

###### Lemma 2.1.

Let $`n`$, $`1kn`$. Then

$$m_{𝔻^n}(A_{𝔻^n,k},z)=g_{𝔻^n}(A_{𝔻^n,k},z)=|z_1|\mathrm{}|z_k|,z𝔻^n.$$

###### Proof of Lemma 2.1.

(Due to P.Pflug.) obviously,

$$|z_1|\mathrm{}|z_k|m_{𝔻^n}(A_{𝔻^n,k},z)g_{𝔻^n}(A_{𝔻^n,k},z).$$

It remains to prove that $`g_{𝔻^n}(A_{𝔻^n,k},z)|z_1|\mathrm{}|z_k|`$.

It is enough to prove that for any $`1kn`$ and for any function $`u:𝔻^n[0,1)`$ such that $`\mathrm{log}u𝒫𝒮(𝔻^n)`$ and $`u(z)C(a)za`$, $`aA_{𝔻^n,k}`$, $`z𝔻^n`$, we have $`u(z)|z_1|\mathrm{}|z_k|`$. We proceed by induction on $`k`$.

For $`k=1`$ the inequality follows from the Schwarz type lemma for logarithmically subharmonic functions $`u(,z_2,\mathrm{},z_n)`$, $`z_2,\mathrm{},z_n𝔻`$.

For $`k>1`$ we first apply the case $`k=1`$ and get $`u(z_1,\mathrm{},z_n)|z_1|`$, $`z𝔻^n`$. Applying the inductive assumption to $`u(z_1,)/|z_1|`$, $`z_1𝔻`$, finishes the proof. ∎

Take a $`u:𝔼[0,1)`$ with $`\mathrm{log}u𝒫𝒮(𝔼)`$ and $`u(\zeta )C(a)\zeta a`$, $`aA`$, $`\zeta 𝔼`$. Consider the mapping

$$\begin{array}{c}\iota _z:𝔻^d(\zeta _1,\mathrm{},\zeta _d)\hfill \\ \hfill (\zeta _1\left(\frac{c_d(z)}{2p_1}\right)^{\frac{1}{2p_1}},\mathrm{},\zeta _d\left(\frac{c_d(z)}{2p_d}\right)^{\frac{1}{2p_d}},z_{d+1},\mathrm{},z_n)𝔼.\end{array}$$

Applying the holomorphic contractivity of the generalized Green function and Lemma 2.1 proves that $`m_𝔼(A,z)g_𝔼(A,z)R_𝔼(A,z)`$.

It remains to prove that $`m_𝔼(A,z)R_𝔼(A,z)`$ (resp. $`g_𝔼(A,z)R_𝔼(A,z)`$). In the case $`d=k=n`$ it suffices to take

$`f(\zeta ):={\displaystyle \underset{j=1}{\overset{d}{}}}\zeta _j\left({\displaystyle \frac{2p_j}{c_d(z)}}\right)^{\frac{1}{2p_j}}.`$
Then $`f𝒪(𝔼,𝔻)`$ and $`|f(z)|=R_𝔼(A,z)`$.

It remains to prove Theorem 1.2 in the remaining cases. We may assume that $`z_j0`$, $`j=1,\mathrm{},d`$. Consider the following lemma.

###### Lemma 2.2.

Define $`𝔼^{}:=𝔼_{p_{d+1},\mathrm{},p_n}`$.

1. Let $`v:𝔼^{}[0,1)`$, $`v0`$ be such that $`\mathrm{log}v𝒫𝒮(𝔼^{})`$ and $`v(\zeta )|\zeta _j|`$, $`\zeta 𝔼^{}`$, $`j=d+1,\mathrm{},k`$, and the mapping

$$𝔼^{}(\zeta _{d+1},\mathrm{},\zeta _n)v(\zeta _{d+1},\mathrm{},\zeta _n)r_d^{q_d}(\zeta )$$

attains its maximum value $`M`$ at $`(z_{d+1},\mathrm{},z_n)`$. Then $`g_𝔼(A,z)=R_𝔼(A,z)`$.
2. The following conditions are equivalent:

* There exists an $`h𝒪(𝔼^{})`$, $`h0`$ such that $`h(\zeta )=0`$ for $`\zeta _{d+1}\mathrm{}\zeta _k=0`$ and the mapping

    $$𝔼^{}(\zeta _{d+1},\mathrm{},\zeta _n)|h(\zeta _{d+1},\mathrm{},\zeta _n)|r_d^{q_d}(\zeta )$$

    attains its maximum value $`M`$ at the point $`(z_{d+1},\mathrm{},z_n)`$,
    * $`m_𝔼(A,z)=R_𝔼(A,z)`$.

We present a proof in Section 3. The result above reduces the proof of Theorem 1.2 to the following propositions.

###### Proposition 2.3.

There exists a $`v`$ as required in Lemma 2.2(a).

###### Proposition 2.4.

Using notation of Theorem 1.2, assume that $`p_j1/2`$, $`j=d+1,\mathrm{},n`$. Then there exists an $`h𝒪(𝔼^{})`$ as required in Lemma 2.2(b).

###### Proposition 2.5.

Using notation of Theorem 1.2, assume that $`p_{k+1}<1/2`$. Then one cannot find an $`h`$ as required in Lemma 2.2(b) for $`|z_{\mathrm{}}|0`$ small enough, $`\mathrm{}=1,\mathrm{},k+1`$ and $`z_{\mathrm{}}=0`$, $`\mathrm{}=k+2,\mathrm{},n`$.

###### Proposition 2.6.

Using notation of Theorem 1.2, assume that $`n=k=2`$. Additionally, assume that $`p_21/2`$ or $`8p_1+4p_2(1p_2)>1`$. Then there exists an $`h𝒪(𝔻)`$ as required in Lemma 2.2(b).

## 3. Proof of Lemma 2.2

(a) Put

$$u(\zeta _1,\mathrm{},\zeta _n):=M^1\left(\underset{j=1}{\overset{d}{}}(2p_j)^{\frac{1}{2p_j}}|\zeta _j|\right)q_d^{q_d}v(\zeta _{d+1},\mathrm{},\zeta _n).$$

Obviously $`\mathrm{log}u𝒫𝒮(𝔼)`$ and $`u(\zeta )C|\zeta _j|C\zeta a`$, $`\zeta 𝔼`$, whenever $`a_j=0`$ for some $`j\{1,\mathrm{},k\}`$. For $`\zeta 𝔼`$ we have:

$$u(\zeta )M^1\left(\frac{_{j=1}^d|\zeta _j|^{2p_j}}{q_d}\right)^{q_d}q_d^{q_d}v(\zeta _{d+1},\mathrm{},\zeta _n)<1.$$

Consequently $`u:𝔼[0,1)`$. On the other hand:

$$u(z)=M^1R_𝔼(A,z)r_d^{q_d}(z)v(z_{d+1},\mathrm{},z_n)=R_𝔼(A,z).$$

(b) Assume that such an $`h`$ exists. Put

$$f(\zeta _1,\mathrm{},\zeta _n):=M^1\left(\underset{j=1}{\overset{d}{}}(2p_j)^{\frac{1}{2p_j}}\zeta _j\right)q_d^{q_d}h(\zeta _{d+1},\mathrm{},\zeta _n).$$

Observe that $`f(\zeta )=0`$ for $`\zeta A`$. Similarly as in (a) we prove that $`|f(\zeta )|<1`$, $`\zeta 𝔼`$ and $`f(z)=R_𝔼(A,z)`$.

Assume now that $`m_𝔼(A,z)=R_𝔼(A,z)`$. Let $`f𝒪(𝔼,𝔻)`$ be such that $`f|_A0`$ and $`|f(z)|=R_𝔼(A,z)`$ (cf. \[Jar-Jar-Pfl 2003\], Property 2.5). Put

$$h(\zeta ):=\frac{^df}{z_1\mathrm{}z_d}(0,\zeta ).$$

By definition, we have $`h(\zeta )=0`$ for $`\zeta _{d+1}\mathrm{}\zeta _k=0`$. Applying the Schwarz lemma to the mapping $`f\iota _\zeta `$, $`\zeta 𝔼^{}`$, we get

$`|h(\zeta _{d+1},\mathrm{},\zeta _n)|r_d^{q_d}(\zeta )`$ $`\left({\displaystyle \underset{j=1}{\overset{d}{}}}(2p_j)^{\frac{1}{2p_j}}\right)q_d^{q_d},\zeta 𝔼^{},`$
$`|h(z_{d+1},\mathrm{},z_n)|r_d^{q_d}(z)=`$ $`\left({\displaystyle \underset{j=1}{\overset{d}{}}}(2p_j)^{\frac{1}{2p_j}}\right)q_d^{q_d}.`$
## 4. Proof of Propositions 2.3, 2.4, 2.5, and 2.6

###### Proof of Proposition 2.3.

We may assume that $`z_{d+1},\mathrm{},z_n0`$. Consider functions of the form

$$v_\alpha (\zeta _{d+1},\mathrm{},\zeta _n)=\left(\underset{j=d+1}{\overset{k}{}}|\zeta _j|^{1+\alpha _j}\right)\left(\underset{j=k+1}{\overset{n}{}}|\zeta _j|^{\alpha _j}\right),$$

where $`\alpha =(\alpha _{d+1},\mathrm{},\alpha _n)`$, $`\alpha _{d+1},\mathrm{},\alpha _n0`$. Obviously $`v:𝔼^{}[0,1)`$, $`\mathrm{log}v𝒫𝒮(𝔼^{})`$, and $`v(\zeta )|\zeta _j|\zeta a`$, $`\zeta 𝔼^{}`$, whenever $`a_j=0`$ for some $`j\{d+1,\mathrm{},k\}`$. Since $`v_\alpha (\zeta _{d+1},\mathrm{},\zeta _n)=v_\alpha (|\zeta _{d+1}|,\mathrm{},|\zeta _n|)`$, it is enough to find an $`\alpha `$ such that the function

$$𝔼^{}_+^{nd}(t_{d+1},\mathrm{},t_n)v_\alpha (t_{d+1},\mathrm{},t_n)r_d^{q_d}(t)$$

attains its maximum at $`(z_{d+1},\mathrm{},z_n)`$. Considering the partial (logarithmic) derivatives results in the following equations:

$`0=1+`$ $`\alpha _j2p_jq_d{\displaystyle \frac{t_j^{2p_j}}{r_d(t)}},j=d+1,\mathrm{},k,`$
$`0=`$ $`\alpha _j2p_jq_d{\displaystyle \frac{t_j^{2p_j}}{r_d(t)}},j=k+1,\mathrm{},n.`$
These give formulas for $`\alpha _{d+1},\mathrm{},\alpha _n`$ such that $`(z_{d+1},\mathrm{},z_n)`$ is the common zero of the derivatives. To prove that there are no other points like this, consider a reformulation of the above equations:

$`r_d(t)`$ $`={\displaystyle \frac{2p_jq_dt_j^{2p_j}}{1+\alpha _j}},j=d+1,\mathrm{},k,`$
$`r_d(t)`$ $`={\displaystyle \frac{2p_jq_dt_j^{2p_j}}{\alpha _j}},j=k+1,\mathrm{},n.`$
The left side is decreasing in any of the variables, while the right sides are increasing. Thus, at most one common zero is allowed.

It remains to check, whether $`\alpha _j0`$, $`j=d+1,\mathrm{},n`$. Obviously, this is true for $`j=k+1,\mathrm{},n`$ and in the remaining cases we have:

$$\alpha _j=\frac{2p_jq_dt_j^{2p_j}r_d(t)}{t_jr_d(t)}0,$$

since this is the way we have chosen $`d`$. ∎

###### Proof of Proposition 2.4.

We may assume that $`z_{d+1},\mathrm{},z_n0`$. Consider functions of the form

$$h_\alpha (\zeta _{d+1},\mathrm{},\zeta _n)=\left(\underset{j=d+1}{\overset{k}{}}\zeta _je^{\alpha _j\zeta _j}\right)\left(\underset{j=k+1}{\overset{n}{}}e^{\alpha _j\zeta _j}\right),$$

where $`\alpha =(\alpha _{d+1},\mathrm{},\alpha _n)`$, $`\alpha _{d+1},\mathrm{},\alpha _n0`$. Since $`|h_\alpha (\zeta _{d+1},\mathrm{},\zeta _n)|h_\alpha (|\zeta _{d+1}|,\mathrm{},|\zeta _n|)`$, it is enough to find an $`\alpha `$ such that

$$𝔼^{}_+^{nd}(t_{d+1},\mathrm{},t_n)h_\alpha (t_{d+1},\mathrm{},t_n)r_d^{q_d}(t)$$

attains its maximum at $`(z_{d+1},\mathrm{},z_n)`$. Considering the partial (logarithmic) derivatives results in the following equations:

$`0={\displaystyle \frac{1}{t}}_j+`$ $`\alpha _j2p_jq_d{\displaystyle \frac{t_j^{2p_j1}}{r_d(t)}},j=d+1,\mathrm{},k,`$
$`0=`$ $`\alpha _j2p_jq_d{\displaystyle \frac{t_j^{2p_j1}}{r_d(t)}},j=k+1,\mathrm{},n.`$
We continue as in the proof of Proposition 2.3. ∎

###### Proof of Proposition 2.5.

Consider the following two lemmas.

###### Lemma 4.1.

Assume that there exist $`0<c<b<1`$ such that the function

$$\phi :[0,b]t(1t^p)^q$$

is strictly concave and

$$\phi (0)+\frac{b}{c}(\phi (c)\phi (0))>\phi (b)+2.$$

Let $`f𝒪(𝔻)`$, $`f0`$ such that $`|f(\zeta )|/\phi (|\zeta |)`$ attains its maximum at $`w_0`$. Then $`w_0=0`$ or $`|w_0|c`$.

###### Proof.

Assume that $`|w_0|(0,c)`$. We may assume that $`|f(w_0)|=\phi (|w_0|)`$. Consider the function

$$\psi :[0,b]t|f(0)|+\frac{t}{|w_0|}|f(w_0)f(0)|.$$

From $`\psi (0)\phi (0)`$, $`\psi (|w_0|)\phi (|w_0|)`$, and the convexity condition we get:

$`\psi (b)=`$ $`|f(0)|+{\displaystyle \frac{b}{|w_0|}}|f(w_0)f(0)|`$
$``$ $`\phi (0)+{\displaystyle \frac{b}{|w_0|}}|\phi (|w_0|)\phi (0)|`$
$``$ $`\phi (0)+{\displaystyle \frac{b}{c}}|\phi (c)\phi (0)|>\phi (b)+2.`$
The Schwarz lemma and the maximum principle imply that there exists a $`w𝔻`$ with $`|w|=b`$ and

$$\frac{|f(w)f(0)|}{|w|}\frac{|f(w_0)f(0)|}{|w_0|}.$$

This means that

$`|f(w)|`$ $`|f(w)f(0)||f(0)|=|f(0)|+|f(w)f(0)|2|f(0)|`$
$``$ $`\psi (b)2|f(0)|>\phi (b)+22|f(0)|\phi (b)=\phi (|w|).`$
This contradicts the maximality of $`|f(\zeta )|/\phi (|\zeta |)`$ at $`w_0`$. ∎

###### Lemma 4.2.

Assume that $`p(0,1)`$, $`q>0`$. Then there exist $`b(0,1)`$, $`c(0,b)`$ as required in Lemma 4.1.

###### Proof.

We have

$$\phi ^{}(t)=\frac{pq}{t^{1p}(1t^p)^{q+1}}$$

which is decreasing for small $`t`$. To prove the existence of $`c`$ observe that

$$\underset{d0^+}{lim}\frac{\phi (d)\phi (0)}{d}=+\mathrm{}.$$

Observe that $`d=k`$ for $`|z_{\mathrm{}}|`$ small enough, $`\mathrm{}=1,\mathrm{},k`$. Let $`h𝒪(𝔼^{})`$ and consider $`f(\zeta ):=h(\zeta ,0,\mathrm{},0)`$, $`\zeta 𝔻`$. Put $`p:=2p_{k+1}`$, $`q=q_d`$. Let $`b,c`$ be as in Lemma 4.2. It follows from Lemma 4.1 that $`0|w_0|<c`$ cannot be a maximum of $`𝔻\zeta |f(\zeta )|/\phi (|\zeta |)`$. In particular $`z`$ cannot be a maximum of $`𝔼^{}\zeta |h(\zeta )|r_d^{q_d}(\zeta )`$ for $`0|z_{k+1}|<c`$. ∎

###### Proof of Proposition 2.6.

Consider the following lemma.

###### Lemma 4.3.

Let $`a,c>0`$, $`t_0(0,1)`$ be such that $`c1`$ or $`4a+2c>1+c^2`$, $`t_0^c>\tau :=a/(a+c)`$. Then there exist $`b>0`$ and $`r1`$ such that

$$h:[0,1]t\frac{t^a}{(rt)^b}(1t^c)[0,1]$$

admits its maximum at $`t_0`$.

###### Proof.

Comparing $`\phi :=th^{}(t)/h(t)`$ to zero we get the equation:

$$\phi (r,t)=a+\frac{bt}{rt}\frac{ct^c}{1t^c}=0.$$

This gives us a formula for $`r`$:

$$r(t)=\frac{t(bbt^ca+at^c+t^cc)}{a+at^c+t^cc}.$$

Observe that:

$`\underset{t\tau ^{1/c}}{lim}r(t)`$ $`=+\mathrm{},`$
$`\underset{t1}{lim}r(t)`$ $`=1.`$
In order to prove that $`r(t)>1`$ and that $`\phi (r,)`$ has only one zero it suffices to show that $`r^{}(t)<0`$. We have:

$$r^{}(t)=\frac{((a+c)t^{2c}+(2a+cc^2)t^ca)b+(at^c+t^cca)^2}{(at^c+t^cca)^2}.$$

It remains to show that the coefficient $`\alpha (t^c)`$ next to $`b`$ is negative. We have:

$`\alpha (\tau )`$ $`={\displaystyle \frac{ac^2}{a+c}}<0,`$
$`\alpha (1)`$ $`=c^2<0,`$
$`\alpha ^{}(u)`$ $`=2(a+c)u+2a+cc^2.`$
Let $`u_0`$ be the zero of $`\alpha ^{}(u)`$. For $`c1`$ we have $`u_0\tau `$ and we are done. Otherwise $`u_0(\tau ,1)`$ and $`4(a+c)\alpha (u_0)=c^2(1+c^24a2c)<0`$. ∎

We may assume that $`z_0>0`$. Put $`a=2p_1`$, $`c=2p_2`$, $`t_0=z_0`$. Let $`r`$ be as in Lemma 4.3. Putting

$$h(\zeta )=\frac{\zeta }{(r\zeta )^{b/a}}.$$

completes the proof. ∎

###### Remark 4.4.

The calculations in the proof of Proposition 2.6 can be performed using alternative function families e.g.

$`h_j(\zeta ):=`$ $`(\zeta +r)^k,r0,k=0,1,2,\mathrm{},`$
$`h_j(\zeta ):=`$ $`(\zeta +1)^\alpha ,\alpha >0,`$
$`h_j(\zeta ):=`$ $`(r\zeta )^\alpha \zeta ^k,\alpha <0,r1,k=0,1,2,\mathrm{}`$
$`h_j(\zeta ):=`$ $`\left({\displaystyle \frac{\zeta +\delta }{1+\delta \zeta }}\right)\zeta ^k,\delta [0,1],k=0,1,2,\mathrm{}.`$
However, the author was unable to solve the general case using any of them.

Acknowledgement. The author is indebted to professor Peter Pflug for his valuable comments on the paper.


text:
# Random walks in random environment: What a single trajectory tells

## 1. Introduction

Suppose you are given a “typical” trajectory of a random walk in an iid random environment. Can you say what the law of the environment is on the basis of the information supplied by this single trajectory? Can you determine the law of the walk? Such questions may arise if one intends to use the random environment model in applications.

These questions are essentially pointless if the group is finite (in which case the environment at each of the finitely many sites that happen to be visited infinitely many times can of course be determined, but it is hard to say much more). So we assume that the group is infinite, and we go a little further: we assume that the (random) set of sites visited by the walk is almost surely infinite. (See remark 5.1.)

Questions of this kind have been studied in the context of random walks in random scenery by Benjamini and Kesten , Löwe and Matzinger , and Matzinger .

In the case of an iid random environment, the information furnished by a single (“typical”) trajectory tells us whether the walk is recurrent; indeed, one can show that one of the events {each visited site is visited infinitely many times}, {no site is visited infinitely many times} is an almost sure event. (Cf Kalikow .)

Now, if the walk is recurrent, the problem is quite simple: we can know much more than the *law* of the environment, because we find the environment *itself* at each visited site, which is given by the frequency of each possible jump from this site. In the transient case, the “naïve” approach consisting of doing statistics on sites which have been visited many times cannot be utilized directly, since the assumption of being at a site which has been frequently visited introduces a bias on the environment at that site, which should encourage jumps to sites from which it is be easier to come back (loosely speaking, close sites).

We present a procedure that eliminates any source of bias, collecting information on sites displaying some specified “histories”. Each such “history” which $`can`$ be encountered $`is`$ almost surely encountered infinitely often (Proposition 3). This is combined with an interpretation of the process as an edge-oriented reinforced random walk (cf Enriquez and Sabot ), allowing us to find the exact law of the process. Now, there may exist “bad” transitions: if the walk jumps from a site along a “bad” transition, it will never get back to that site. If the set of these “bad” transitions is empty or has just one element, we can find the distribution of the environment (Theorem 1).

Finally, we show how countably many independent trajectories having the distribution of the original walk can be generated by concatenating steps of the observed trajectory. The algorithm is purely “mechanical”: it does not imply any computation, and, in particular, the knowledge of the law of the walk (or of the environment) is not needed.

## 2. Framework and notations

Our “canonical” process $`X:=(X_n)_{n0}`$ walks on a group $`G`$. We denote by $`(_n)_{n0}`$ the natural filtration of $`X`$.

We assume that the group $`G`$ is Abelian, although this is never used in our arguments. Its only utility is in the possibility of writing things like $`xy=e`$ or $`x=y+e`$ indifferently.

We use the additive notation, and the identity element of $`G`$ is denoted by $`0`$.

We assume moreover that the group $`G`$ is countable. This assumption can be dispensed with$``$see remark 5.2$``$but we feel that it renders the reading easier and the discussion more tractable. It does not affect the core of the argument.

We denote by $``$ the set $`\{1,2,\mathrm{}\}`$.

### 2.1. Random walks in random environment.

We denote by $`𝒫`$ the set of non-negative families $`p:=(p_e)_{eG}`$ such that $`_{eG}p_e=1`$. The environment at the site $`x`$, $`\nu (x):=(\nu (x,e))_{eG}`$, is a random element of $`𝒫`$. We assume that the environments at sites are iid $`𝒫`$-valued random variables with common distribution $`\mu `$.

We let $`\nu _e:=\nu (0,e)`$.

The random environment $`\nu :=(\nu (x))_{xG}`$ is a random element of $`𝒫^G`$, and it is governed by the probability measure $`\mu ^G`$.

For all $`\pi =((\pi (x,e))_{eG})_{xG}𝒫^G`$, let $`Q_\pi `$ be the probability measure under which $`X`$ is a $`G`$-valued Markov chain started at 0 whose transition probability from $`x`$ to $`x+e`$ is $`\pi (x,e)`$ ($`x,eG`$).

The law of the random walk in random environment (or the so-called “*annealed*” law) is the probability measure $`𝐏^\mu =Q_\pi \mu ^G(d\pi )(=𝐄[Q_\nu ]`$$`Q_\nu `$ being what is usually called the “quenched” law).

We recall our “infinitude assumption” according to which the (random) set of sites visited by the walk, $`S:=\{X_n|n0\}`$, is $`𝐏^\mu `$-almost surely infinite. (This implies, of course, that the group $`G`$ itself is infinite.)

We let $`E`$ denote the set of those $`gG`$ such that the probability of the event $`\{\nu _g>0\}`$ is strictly positive. (It is easy to see that the random set $`\{X_{n+1}X_n|n0\}`$ is $`𝐏^\mu `$-almost surely exactly $`E`$.) We then partition $`E`$ into two sets, $`R`$ and $`T`$, defined as follows.

$``$ $`R`$ is the set of elements $`r`$ of $`E`$ that can be written as $`(e_1+\mathrm{}+e_n)`$ where $`(e_i)_{1in}`$ is a finite nonempty sequence of elements of $`E`$. It is easy to see that $`rR`$ if and only if $`𝐏^\mu (X_1=r\text{ and, for some }n>1,X_n=0)>0`$ ; and Proposition 3 below implies that if $`rR`$, then the random set $`\{n|X_{n+1}=X_n+r\text{ and, for some }k>0,X_{n+k}=X_n\}`$ is $`𝐏^\mu `$-almost surely infinite. It is therefore quite easy to identify $`R`$ when observing a single trajectory.

$``$ $`T`$ is the complement of $`R`$ in $`E`$. It represents the “possible” transitions which do not allow a return to the original site.

### 2.2. Histories.

We start with some definitions.

###### Definition 1.

The history of the site $`x`$ at time $`n`$, which we denote by $`H(n,x)`$, is the random finite sequence of elements of $`G`$ defined by the successive moves of the process from the site $`x`$ before time $`n`$. More formally, $`H(0,x)`$ is the empty sequence $`()`$; $`H(n+1,x)=H(n,x)`$ if $`X_nx`$; and, if $`X_n=x`$, then $`H(n+1,x)`$ is the finite sequence obtained by adjoining $`X_{n+1}X_n`$ as a new rightmost term to $`H(n,x)`$.

Let us denote by $`(_+^G)_0`$ the set of families $`(n_g)_{gG}_+^G`$ with a finite number of non null terms.

###### Definition 2.

The unordered history of the site $`x`$ at time $`n`$, denoted by $`\stackrel{}{N}(n,x):=(N_g(n,x))_{gG}`$, is a random element of $`(_+^G)_0`$ where, for all $`gG`$, $`N_g(n,x)`$ is the random number of moves from the site $`x`$ to $`x+g`$ before time $`n`$. In other words, $`N_g(n,x)=_{l=0}^{n1}1_{\{X_l=x,X_{l+1}X_l=g\}}`$.

Also, the local unordered history at time $`n`$ is the unordered history of the site $`X_n`$ at time $`n`$, $`\stackrel{}{N}(n):=\stackrel{}{N}(n,X_n)`$.

### 2.3. Reinforced random walks.

An edge-oriented reinforced random walk consists of a discrete random process whose transition probabilities are functions of the number of each type of move in the history of the process that has been made from the site currently occupied. A good point of view, in order to get a non-biased procedure of reconstitution of the environment, is to view the random walk in random environment as an edge-oriented reinforced random walk. It is the essence of the easy part of the result of  (the other part examines the conditions on a reinforced random walk to correspond to a RWRE).

We introduce the reinforced random walks by the following definitions.

###### Definition 3.

A reinforcement function is a function

$$V:(_+^G)_0𝒫$$

$$\stackrel{}{n}=(n_g)_{gG}V(\stackrel{}{n}):=(V_e(\stackrel{}{n}))_{eG}$$

###### Definition 4.

We call edge-oriented reinforced random walk with reinforcement function $`V`$ the random walk defined by the law $`𝐏^V`$ on the trajectories starting at $`0`$ given by

$$𝐏^V(X_{n+1}X_n=e|_n)=V_e(\stackrel{}{N}(n)).$$

## 3. Tools

### 3.1. RWRE as an edge-oriented reinforced random walk.

We can now state the result of Enriquez and Sabot :

###### Proposition 1.

The annealed law $`𝐏^\mu `$ of the RWRE coincides with the law $`𝐏^V`$ of the reinforced random walk whose reinforcement function $`V`$ satisfies, for all $`eG`$,

$$V_e(\stackrel{}{n})=\frac{𝐄[\nu _e\underset{gG}{}\nu _g^{n_g}]}{𝐄[_{gG}\nu _g^{n_g}]}$$

whenever $`\stackrel{}{n}(_+^G)_0`$ such that $`𝐄[_{gG}\nu _g^{n_g}]>0`$.

In order to be self-contained we recall the proof of this proposition.

PROOF.  For any $`x,e`$ in $`G`$, for all $`n,`$ $`𝐏^\mu `$-almost everywhere on the event $`\{X_n=x\}`$,

$$𝐏^\mu (X_{n+1}=x+e|_n)=\frac{𝐄[\nu (x,e)\underset{yG}{}\underset{gG}{}\nu (y,g)^{n_g(n,y)}]}{𝐄[_{yG}_{gG}\nu (y,g)^{n_g(n,y)}]}$$

Now using the independence of the random variables $`\nu (y)`$ for different sites $`y`$, the terms depending on $`\nu (y)`$ for $`yx`$ cancel in the previous ratio, and we get the result. ∎

The following result is an analogue of the strong Markov property for reinforced random walks.

###### Proposition 2.

Let $`X`$ be a reinforced random walk with reinforcement function $`V`$, and let $`T`$ be a stopping time with respect to the natural filtration of $`X`$. Assume $`T`$ is almost surely finite. Then

$$𝐏^V(X_{T+1}X_T=e|_T)=V_e(\stackrel{}{N}(T))𝐏^Va.s.$$

The proof is obtained in an obvious way, by considering the events $`\{T=n\}`$.

### 3.2. A zero-one result.

The following zero-one result happens to be quite useful.

###### Proposition 3.

Let $`(r_1,\mathrm{},r_l)`$ a finite (eventually empty) sequence of elements of $`R`$. Let $`S_{(r_1,\mathrm{},r_l)}`$ be the random set $`\{xG|n0,H(n,x)=(r_1,\mathrm{},r_l)\}`$. Then $`S_{(r_1,\mathrm{},r_l)}`$ is either $`𝐏^\mu `$-almost surely empty or $`𝐏^\mu `$-almost surely infinite.

PROOF.  Suppose that $`S_{(r_1,\mathrm{},r_l)}`$ is not $`𝐏^\mu `$-almost surely empty.

This implies that there exists a list of transitions

$$L:=(r_1,e_{1,1},e_{1,2},\mathrm{},e_{1,k_1},r_2,e_{2,1}\mathrm{},e_{2,k_2},r_3,\mathrm{}..,e_{l1,k_{l1}},r_l)$$

such that, for all $`m\{1,\mathrm{},l1\}`$,

$$r_m+\underset{i=1}{\overset{k_m}{}}e_{m,i}=0\text{and, for all}k\{1,\mathrm{},k_m1\},r_m+\underset{i=1}{\overset{k}{}}e_{m,i}0$$

and

$$\gamma :=𝐄[\underset{k=1}{\overset{l}{}}\nu _{r_k}]𝐄[\underset{m=1}{\overset{l1}{}}\underset{i=1}{\overset{k_m1}{}}\nu (r_m+e_{m,1}+\mathrm{}+e_{m,i},e_{m,i+1})]>0.$$

(Note that if $`r_m=0`$, then $`k_m=0`$.)

Let $`q:=l+k_1+\mathrm{}+k_{l1}`$ be the length of the list $`L`$.

For all $`k\{1,\mathrm{},q\}`$, denote by $`y_k`$ the $`k`$-th term of the list $`L`$. For all $`k\{1,\mathrm{},q+1\}`$, set $`x_k:=_{i=1}^{k1}y_i`$ . ($`x_1=0`$.)

Now consider the list $`(g_1,g_2,\mathrm{})`$ of newly visited sites in their order of appearance. So $`𝒮=\{g_1,g_2,\mathrm{}\}`$. By the assumption made in the introduction, $`𝒮`$ is almost surely infinite.

To any integer $`n1`$, we associate a random integer $`i(n)`$ defined by

$$i(n):=\mathrm{min}\{i1|k\{1,\mathrm{},q\},g_i=g_n+x_k\}.$$

We denote by $`k(n)`$ the random smallest integer $`m1`$ such that $`g_{i(n)}=g_n+x_m`$. (Obviously, $`k(n)q`$.) The sequence $`(g_{i(n)})_{n1}`$ takes infinitely many values (since the infinite set $`𝒮`$ is included in $`\{g_{i(1)},g_{i(2)},\mathrm{}\}\{x_1,\mathrm{},x_q\}`$). Now, for any $`i1`$, we denote by $`T_i`$ the hitting time of $`g_i`$ by the walk. By definition of $`i(n)`$, none of the sites $`g_{i(n)}x_{k(n)}+x_j(1jq)`$ is visited by the trajectory before time $`T_{i(n)}`$.

As a result, there exist infinitely many sites $`g_1^{},g_2^{},\mathrm{}`$ visited by the trajectory (enumerated in their order of appearance) such that, for some $`k\{1,\mathrm{},q\}`$, if $`T_n^{}`$ denotes the hitting time of $`g_n^{}`$, then none of the sites $`g_n^{}x_k+x_j(1jq)`$ is visited by the trajectory before time $`T_n^{}`$. We denote by $`k^{}(n)`$ the least such integer $`k`$. $`T_n^{}`$ are clearly stopping times.

For all $`n1`$, let $`\psi _n`$ be the Bernoulli variable that equals 1 if and only if

$$X_{T_n^{}+i}=\{\begin{array}{cc}g_n^{}x_{k^{}(n)}+x_{i+k^{}(n)}\hfill & ifi\{1,\mathrm{},qk^{}(n)\},\hfill \\ g_n^{}x_{k^{}(n)}+x_{ik^{}(n)q}\hfill & ifi\{q+1k^{}(n),\mathrm{},2qk^{}(n)+1\}.\hfill \end{array}$$

(Otherwise, $`\psi _n`$ equals 0.)

Observe that for all $`n`$, if $`\psi _n=1`$, then $`X_{T_n^{}}x_{k^{}(n)}S_{(r_1,\mathrm{},r_l)}`$.

Due to the fact that the prescribed path the process has to follow during the time period $`[T_n^{},T_n^{}+2qk^{}(n)+1]`$ in order to satisfy $`\psi _n=1`$ is a path that does not intersect the trajectory before $`T_n^{}`$,

$$𝐏(\psi _n=1|_{T_n^{}})𝐄[\underset{k=1}{\overset{q}{}}(\nu (x_k,y_k))^2]𝐄[\underset{k=1}{\overset{q}{}}\nu (x_k,y_k)]^2=\gamma ^2.$$

Let $`\xi _n:=\psi _{(2q+1)n}`$ and $`\tau _n=T_{(2q+1)n}^{}`$ ($`n1`$).

For all $`n`$, $`\xi _1,\mathrm{},\xi _n`$ are measurable with respect to $`_{\tau _{n+1}}`$.

It is now obvious that for all $`n,k1`$

$$𝐏(\xi _{n+1}=\mathrm{}=\xi _{n+k}=0)=$$

$$𝐏(\xi _{n+1}=0)𝐏(\xi _{n+2}=0|\xi _{n+1}=0)\mathrm{}𝐏(\xi _{n+k}=0|\xi _{n+1}=\mathrm{}=\xi _{n+k1}=0)(1\gamma ^2)^k.$$

Therefore, almost surely, infinitely many of the $`\psi _n`$’s are equal to 1. This implies that $`S_{(r_1,\mathrm{},r_l)}`$ is almost surely infinite. ∎

REMARK.  If $`G`$ equals $`^d`$, the notion of convexity can be exploited in a proof slightly different from the one given above.

We deduce that the sets $`R`$ and $`T`$ can be “viewed” on the trajectory:

###### Corollary 1.

The sets $`R`$ and $`T`$ are such that

$`R\stackrel{a.s}{=}\{gE|S_{(g)}\text{ is infinite}\}\stackrel{a.s}{=}\{gE|S_{(g)}\mathrm{}\}`$ and $`T\stackrel{a.s}{=}\{gE|S_{(g)}=\mathrm{}\}`$.

Let $`S_\stackrel{}{n}`$ denote the random set $`\{xG|n0,\stackrel{}{N}(n,x)=\stackrel{}{n}\}`$. An easy corollary of the above proposition is the following analogous result concerning unordered histories.

Proposition 3’. If $`\stackrel{}{n}(_+^G)_0`$, then $`S_\stackrel{}{n}`$ is either $`𝐏^\mu `$-almost surely empty or $`𝐏^\mu `$-almost surely infinite.

We now distinguish a particular subset of $`(_+^G)_0`$,

$$𝒩_{poss}:=\{\stackrel{}{n}(_+^G)_0|S_\stackrel{}{n}\mathrm{}𝐏^\mu a.s.\}.$$

Loosely speaking, $`𝒩_{poss}`$ is the set of “possible” unordered histories for sites that are presently occupied. An element $`\stackrel{}{n}=(n_g)_{gG}`$ of $`(_+^G)_0`$ belongs to $`𝒩_{poss}`$ if $`n_g=0`$ whenever $`gR`$ and if, moreover, it satisfies (any one of) the three following equivalent conditions:

(a) $`𝐏^\mu (n0|\stackrel{}{N}(n,0)=\stackrel{}{n})>0`$;

(b) the random set $`S_\stackrel{}{n}`$ is almost surely infinite;

(c) $`𝐄[_{rR}\nu _r^{n_r}]>0`$.

Note that, by (c) and Proposition 1, the law of the process is determined by the restriction of the reinforcement function to the set $`𝒩_{poss}`$.

## 4. What a single trajectory tells

In the sequel we assume that the law of the process $`X`$ is $`𝐏^\mu `$ (or, equivalently, $`𝐏^V`$).

### 4.1. Determining the law of the walk.

As noted in the introduction, “straightforward” computation based on the frequencies of transitions from sites visited “many” times is not reliable. What we do here, instead, is collecting information on sites displaying some specified histories (or specified *unordered* histories).

For any $`\stackrel{}{n}𝒩_{poss}`$, let $`T_i^\stackrel{}{n}(i1)`$ be the successive times where the unordered history of the currently occupied site is $`\stackrel{}{n}`$:

$``$ $`T_0^\stackrel{}{n}:=inf\{k0|\stackrel{}{N}(k)=\stackrel{}{n}\}`$

$``$ $`i0,T_{i+1}^\stackrel{}{n}:=inf\{k>T_i^\stackrel{}{n}|\stackrel{}{N}(k)=\stackrel{}{n}\}`$
(We ignore the case (happening on the negligible event $`\{S_\stackrel{}{n}`$ is finite}) where some $`T_i^\stackrel{}{n}`$ is infinite.)

###### Proposition 4.

The $`E`$-valued random variables $`\mathrm{\Delta }_i^\stackrel{}{n}:=X_{T_i^\stackrel{}{n}+1}X_{T_i^\stackrel{}{n}}`$ ($`i1`$, $`\stackrel{}{n}𝒩_{poss}`$) are independent. Also, for all $`\stackrel{}{n}𝒩_{poss}`$, the random variables $`\mathrm{\Delta }_i^\stackrel{}{n}`$ have the same law:

$$eE,𝐏(\mathrm{\Delta }_i^\stackrel{}{n}=e)=V_e(\stackrel{}{n}).$$

PROOF.  Let $`\mathrm{\Theta }_i^\stackrel{}{n}`$ ($`i1`$, $`\stackrel{}{n}(_+^G)_0`$) be independent random variables such that for all $`eE`$,

$$𝐏(\mathrm{\Theta }_i^\stackrel{}{n}=e)=V_e(\stackrel{}{n})$$

Now we consider the process $`(Y_n)_{n0}`$ :

$$Y_0=0,Y_{n+1}=Y_n+\mathrm{\Theta }_{\tau (n)}^{\stackrel{}{N}_y(n)},$$

where $`\stackrel{}{N}_y`$ is to the process $`Y`$ what $`\stackrel{}{N}`$ is to the process $`X`$, and

$$\tau (n):=card\{j\{0,\mathrm{},n\}|\stackrel{}{N}_y(j)=\stackrel{}{N}_y(n)\}.$$

$$\begin{array}{cc}\hfill 𝐏(Y_{n+1}=Y_n+e|\sigma (Y_k,kn))& =𝐏(\mathrm{\Theta }_{\tau (n)}^{\stackrel{}{N}_y(n)}=e|\sigma (Y_k,kn))\hfill \\ & =𝐄[_{l0,\stackrel{}{m}(_+^G)_0}1_{\tau (n)=l,\stackrel{}{N}_y(n)=\stackrel{}{m}}1_{\mathrm{\Theta }_l^\stackrel{}{m}=e}|\sigma (Y_k,kn)]\hfill \end{array}$$

But on the event $`A_{l,\stackrel{}{m},n}:=\{\tau (n)=l,\stackrel{}{N}_y(n)=\stackrel{}{m}\}`$, $`\mathrm{\Theta }_l^\stackrel{}{m}`$ is independent of $`\sigma (Y_k,kn)`$. Thus,

$$\begin{array}{cc}\hfill 𝐏(Y_{n+1}=Y_n+e|\sigma (Y_k,kn))& =_{l0,\stackrel{}{m}(_+^G)_0}1_{\tau (n)=l,\stackrel{}{N}_y(n)=\stackrel{}{m}}𝐄[1_{\mathrm{\Theta }_l^\stackrel{}{m}=e}]\hfill \\ & =_{l0,\stackrel{}{m}(_+^G)_0}1_{\tau (n)=l,\stackrel{}{N}_y(n)=\stackrel{}{m}}V_e(\stackrel{}{m})\hfill \\ & =V_e(\stackrel{}{N}_y(n))\hfill \end{array}$$

Consequently, the two processes $`X`$ and $`Y`$ have the law. But $`\mathrm{\Delta }_i^\stackrel{}{n}`$ is to the process $`(X_n)_{n0}`$ exactly what $`\mathrm{\Theta }_i^\stackrel{}{n}`$ is to the process $`Y`$. The result follows. ∎

We deduce from this proposition the following corollary, which describes the construction of the reinforcement function on $`𝒩_{poss}`$ or, equivalently (by Proposition 1), of the annealed law:

###### Corollary 2.

If $`\stackrel{}{n}𝒩_{poss}`$, then almost surely, for all $`eE`$,

$$\frac{1_{\mathrm{\Delta }_1^\stackrel{}{n}=e}+\mathrm{}+1_{\mathrm{\Delta }_m^\stackrel{}{n}=e}}{m}V_e(\stackrel{}{n})asm\mathrm{}.$$

### 4.2. The law of the environment.

The next result follows easily.

###### Theorem 1.

(a) A single trajectory determines almost surely the moments of the form $`𝐄_\mu [\nu _{r_1}^{n_1}\mathrm{}\nu _{r_k}^{n_k}\nu _t^\epsilon ]`$ for all $`r_1,\mathrm{},r_kR`$, $`tT`$, $`n_1,\mathrm{},n_k_+`$, $`\epsilon =0`$ or $`1`$. Moreover, if these moments coincide for two distinct environment distributions, the induced RWRE have the same annealed law (and, consequently, two such environment distributions cannot be distinguished).

(b) If card $`T=0`$ or $`1`$, a single trajectory determines almost surely the distribution of the environment.

PROOF.  (a) By Corollary 2, the restriction of $`V`$ to $`𝒩_{poss}`$ is almost surely determined by a single trajectory. So almost surely, for all $`\stackrel{}{n}=(n_g)_{gG}𝒩_{poss}`$ and for all $`eG`$, a single trajectory determines the moments $`𝐄[(_{gG}\nu _g^{n_g})\nu _e]`$. Moreover, all the other moments of the type $`𝐄_\mu [\nu _{r_1}^{n_1}\mathrm{}\nu _{r_k}^{n_k}\nu _t^\epsilon ]`$ ($`r_jR`$, $`tT`$) are zero. Finally, the restriction of $`V`$ to $`𝒩_{poss}`$ determines the law of the process.

(b) If card $`T=0`$, we get all the moments of the $`\nu _e`$’s. Since these variables are compactly supported variables, this determines all the finite dimensional marginals of the distribution of $`\nu `$.111We recall a standard fact: if $`U_1,\mathrm{},U_l`$ are positive random variables such that $`U_1+\mathrm{}+U_l[0,1]`$ almost surely then, for Lebesgue-almost all $`(a_1,\mathrm{},a_l)^l`$,
$$𝐏(U_1<a_1,\mathrm{},U_l<a_l)=\underset{n\mathrm{}}{lim}\underset{\genfrac{}{}{0pt}{}{k_0,\mathrm{},k_l0}{\genfrac{}{}{0pt}{}{k_0+\mathrm{}+k_l=n}{\frac{k_1}{n}<a_1,\mathrm{},\frac{k_l}{n}<a_l}}}{}\frac{n!}{k_0!\mathrm{}k_l!}𝐄[(1U_1\mathrm{}U_l)^{k_0}U_1^{k_1}\mathrm{}U_l^{k_l}].$$

If card $`T=1`$, we get all the moments involving the $`\nu _r`$’s. And if $`t`$ is the unique element of $`T`$, then $`\nu _t=1_{rR}\nu _r`$, and we get all the moments of $`\nu `$. ∎

When card $`T2`$ the law of the environment can be determined in some cases, but not in general. (Accordingly, Corollary 1 of  should be amended; it holds in fact if card $`T1`$, but not in complete generality.) Here are two examples:

EXAMPLE 1. We consider the two following walks on $``$:

$``$ The first one has a deterministic environment, and moves from $`x`$ with probability $`\frac{1}{2}`$ to $`x+1`$, with probability $`\frac{1}{2}`$ to $`x+2`$.

$``$ The environment in the second walk is coin-tossed independently at each site $`x`$: with probability $`\frac{1}{2}`$, the transition probability to $`x+1`$ is 1, and with probability $`\frac{1}{2}`$, the transition probability to $`x+2`$ is 1.

Here $`T=\{1,2\}`$, and, obviously, the two walks have the same law.

EXAMPLE 2. Again, $`G=`$; for any $`x`$, with probability $`\frac{1}{2}`$, the transition probability from $`x`$ to $`x`$ is equal to $`\frac{1}{2}`$, and the transition probability from $`x`$ to $`x+1`$ is also equal to $`\frac{1}{2}`$; and, with probability $`\frac{1}{2}`$, the transition probability from $`x`$ to $`x+2`$ is equal to 1. In this case, $`T=\{1,2\}`$, but the distribution of the environment is almost surely completely determined by the single trajectory we observe.

We can only see 0-transitions (from a site $`x`$ to itself), 1-transitions ($`xx+1`$) and 2-transitions ($`xx+2`$). So $`\mu `$, which is the law of $`\nu (0)`$, satisfies $`\mu (\nu _0+\nu _1+\nu _2=1)=1`$. The fact that a 0-transition is never followed by a 2-transition tells us that if $`\nu _2>0`$, then $`\nu _0=0`$. Statistics on sites from which there are 0-transitions or 1-transitions reveals that the (conditional) distribution of the number of 0-transitions from such a site is geometric. But a geometric distribution cannot be a nondegenerate convex combination of geometric distributions, and the conditional number of 0-transitions from a visited site $`x`$ for which $`\nu (x,0)`$ is given (and is in $`]0,1[`$) has a geometric distribution. We deduce that there is exactly one value $`a`$ ($`=\frac{1}{2}`$) such that almost surely, for all $`x`$, if $`\nu (x,0)0`$, then $`\nu (x,0)=a`$. And a simple computation shows that the proportion of visited sites from which a 0-transition is possible but does not take place fits exactly with the proportion of visited sites from which the first (and unique) transition is a 1-transition; so $`\mu (\nu _1>0\mathrm{and}\nu _2>0)=0`$. We deduce that $`\mu `$ is a convex combination of $`\mu ^1`$ and $`\mu ^2`$, where $`\mu ^1(\nu _0=\nu _1=\frac{1}{2})=1`$ and $`\mu ^2(\nu _2=1)=1`$. And the coefficients (both equal$`\frac{1}{2}`$) in this convex combination are easily determined.

### 4.3. Sampling iid trajectories.

Here we show that infinitely many independent trajectories drawn under $`𝐏^\mu `$ can be generated by concatenating steps of the single trajectory at our disposal.

For simplicity, we describe the construction of just two trajectories, $`X^1`$ and $`X^2`$. This is enough, since we can do the following: once $`X^1`$ and $`X^2`$ are constructed, leave $`X^1`$ as it is and extract two trajectories, say $`X^3`$ and $`X^4`$, out of $`X^2`$ exactly the way we extracted $`X^1`$ and $`X^2`$ out of $`X`$, then, out of $`X^4`$, get $`X^5`$ and $`X^6`$, and so on; and the family $`(X^1,X^3,X^5,\mathrm{})`$ is exactly what we want.

All we do is construct $`X^1`$ (resp. $`X^2`$) using in the “natural” way the transitions $`\mathrm{\Delta }_i^\stackrel{}{n}`$ with $`i`$ odd (resp. even). More formally, $`X^1`$ and $`X^2`$ are defined as follows:

$$X_0^1=0,X_{n+1}^1=X_n^1+\mathrm{\Delta }_{2\tau ^1(n)1}^{\stackrel{}{N}^1(n)},$$

where $`\stackrel{}{N}^1`$ is to the process $`X^1`$ what $`\stackrel{}{N}`$ is to the process $`X`$, and

$$\tau ^1(n):=\mathrm{card}\{j\{0,\mathrm{},n\}|\stackrel{}{N}^1(j)=\stackrel{}{N}^1(n)\};$$

and, similarly,

$$X_0^2=0,X_{n+1}^2=X_n^2+\mathrm{\Delta }_{2\tau ^2(n)}^{\stackrel{}{N}^2(n)},$$

where $`\stackrel{}{N}^2`$ is to the process $`X^2`$ what $`\stackrel{}{N}`$ is to the process $`X`$, and

$$\tau ^2(n):=\mathrm{card}\{j\{0,\mathrm{},n\}|\stackrel{}{N}^2(j)=\stackrel{}{N}^2(n)\}.$$

The validity of this construction is an immediate consequence of Proposition 4.

## 5. Remarks

### 5.1. Infinitude of $`S`$.

The “infinitude assumption” (according to which the random set $`S`$ of sites visited by $`X`$ is almost surely an infinite set) is made in order to avoid discussing rather trivial cases. (If $`S`$ is finite, then precise knowledge of the environment at *some* sites is almost surely available; but, unless some specific conditions are imposed on $`\mu `$, complete knowledge of $`\mu `$ is out of reach.)

Proceeding along the general lines of the proof of Proposition 3, one can show that the random set $`S`$ is either almost surely infinite or almost surely finite.

### 5.2. Countability of $`G`$.

If we abandon the countability assumption on $`G`$, the set $`E=\{xG|P(X_1=x)>0\}`$ remains countable, and our sampling procedure works just as well. Consequently, $`𝐏^\mu `$ can be determined in non-pathological situations (and in particular if $`G`$ is the real line).222Of course, if $`G`$ is countable, the $`\sigma `$-field we use (without explicitly saying so) is the set of all subsets of $`G`$; and if $`G`$ is the real line, we take the Borel $`\sigma `$-field on the line. But problems may arise if a $`\sigma `$-field on $`G`$ is not specified in advance and there is no “natural” $`\sigma `$-field on $`G`$: the very notion of the law of $`X`$ is problematic (and, in fact, even the notion of the law of $`X_1`$ does not make much sense). But even if a $`\sigma `$-field on $`G`$ is “given”, we aren’t through. What we want is to be able to determine, on the basis of the observation of one realization of a random sequence $`(U_1,U_2,\mathrm{})`$ of iid random variables taking values in $`G`$, the probability distribution of $`U_1`$. Now if the $`\sigma `$-field is generated by some countable $`\pi `$-system of subsets of $`G`$, things are all right. Otherwise, there is no general positive result. (This can also be seen by introducing a new kind of reinforcement function which, to a given unordered history at a site, associates the probability that the next transition falls into some measurable set.) If there is no countable set $`DG`$ such that $`X_1D`$ almost surely, then $`\mu `$ cannot be determined (as one can see after studying the first example of section 4). (If there is some countable set $`DG`$ such that $`X_1D`$ almost surely then, almost surely, each$`X_n`$ is in the subgroup generated by $`D`$; and since this subgroup is countable, all we did is adaptable in an obvious way.)

### 5.3. Structure of $`G`$.

The choice of dealing with RWRE on a *group* captures, we think, the essence of the matter. We could have restricted ourselves to groups like $`^d`$ (or some other subgroups of $`^d`$) without a substantial gain in simplicity. A group structure is suitable (though not absolutely indispensable) if the notion of iid random environment is to make sense. We could have dealt with RWRE on homogenous spaces, or on some trees, without gaining new insight.

### 5.4. Assumptions on the environment.

The requirement that environment at sites are iid can be loosened in various ways.

EXAMPLE.  $`G=`$; there are two laws for the environment at sites, say $`\mu _0`$ and $`\mu _1`$; environments at sites are independent; and $`\nu (n)`$ is governed by $`\mu _0`$ if $`n`$ is even, by $`\mu _1`$ if $`n`$ is odd.

EXAMPLE.  $`G=`$; there are laws $`\mu _0`$, $`\mu _1`$…for the environment at a site; $`K`$ is a random variable taking values in the set $`\{2,3,\mathrm{}\}`$; and, conditioned on $`K`$, the $`\nu (n)`$ are independent and, for all $`n`$, $`\nu (n)`$ is governed by $`\mu _{n(\mathrm{mod}K)}`$.

EXAMPLE.  $`G=`$; the couples $`(\nu (2n),\nu (2n+1))(n)`$ are iid, but $`\nu (0)`$ and $`\nu (1)`$ are *not* independent.

### 5.5. Other reinforcements.

Our results on the determination of the law of the process and on sampling iid trajectories can be extended to various other edge-oriented reinforced random walks that do not correspond to a random environment. Whenever an appropriate analogue of Proposition 3 is valid, things work quite well. (A sufficient condition is strict positivity of the restriction of $`V_r`$ to $`(_+^R)_0`$ for all $`rR`$.)