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# Moduli of suspension spectra

## 1. Introduction

A well-known property of the singular cohomology of a space is that it possesses the structure of a commutative graded ring. There is no corresponding property for the singular cohomology of spectra. This discrepancy between spaces and spectra is connected with the observation that a spectrum $`E`$ can fail to have a “diagonal map” $`EEE`$. Of course, a diagonal map exists when $`E`$ is a suspension spectrum. This motivates

###### Question.

Let $`E`$ be a spectrum. When can we find a based space $`X`$ and a weak equivalence of spectra $`\mathrm{\Sigma }^{\mathrm{}}XE`$? In how many ways?

The first result of this paper gives a criterion for deciding the existence part of the above question in the metastable range. Recall that $`E`$ is $`r`$-connected if its homotopy groups $`\pi _{}(E)`$ vanish when $`r`$. Write $`dimEn`$ if $`E`$ is, up to homotopy, a CW spectrum with cells in dimensions $`n`$. Recall that the $`k`$-th extended power $`D_k(E)`$ is the homotopy orbit spectrum of the symmetric group $`\mathrm{\Sigma }_k`$ acting on the $`k`$-fold smash product $`E^k`$.

###### Theorem A (Existence).

There is an obstruction

$$\delta _E[E,\mathrm{\Sigma }D_2(E)],$$

which is trivial when $`E`$ has the homotopy type of a suspension spectrum.

Conversely, if $`\delta _E=0`$, $`E`$ is $`r`$-connected, $`r1`$ and $`dimE3r+2`$, then $`E`$ has the homotopy type of a suspension spectrum.

Before stating our second result, we comment on the relation between Theorem A and one of the early results of Kuhn (\[Ku, Th. 1.2\]) which says that a connected spectrum $`E`$ is a retract of a suspension spectrum if and only if there is a weak equivalence

$$\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}E\underset{k1}{}D_k(E).$$

Theorem A implies one can remove the word “retract” from Kuhn’s result in the metastable range.

Next consider the collection of pairs $`(X,h)`$ in which $`X`$ is a based space and $`h:\mathrm{\Sigma }^{\mathrm{}}XE`$ is a weak (homotopy) equivalence. Equate two such pairs $`(X,h)`$ and $`(Y,g)`$ if and only if there is a map of spaces $`f:XY`$ such that $`g\mathrm{\Sigma }^{\mathrm{}}f`$ is homotopic to $`h`$ (in particular, $`f`$ is a homology isomorphism). This generates an equivalence relation. Let $`\mathrm{\Theta }_E`$ denote the associated set of equivalence classes.

Our second result identifies $`\mathrm{\Theta }_E`$ in the metastable range.

###### Theorem B (Enumeration).

Assume $`\mathrm{\Theta }_E`$ is nonempty and is equipped with basepoint. Then there is a basepoint preserving function

$$\varphi :\mathrm{\Theta }_E[E,D_2(E)].$$

If $`E`$ is $`r`$-connected, $`r1`$ and $`dimE3r+2`$, then $`\varphi `$ is a surjection. If in addition $`dimE3r+1`$, $`\varphi `$ is a bijection. In particular, $`\mathrm{\Theta }_E`$ has the structure of an abelian group when $`dimE3r+1`$.

Theorem B leads to the possibility of calculating $`\mathrm{\Theta }_E`$ in a number of simple cases. For example, table 4.1 appearing in Mahowald’s memoir \[Mah\] provides extensive calculations of $`\mathrm{\Theta }_E`$ at the prime $`2`$ in the case of two cell spectra (see §8 below).

In fact, $`\mathrm{\Theta }_E`$ is the set of path components of a space $`𝔐_E`$ which is defined as follows: let $`C_E`$ be the category whose objects are pairs $`(Y,h)`$ with $`Y`$ a based space and $`h:\mathrm{\Sigma }^{\mathrm{}}YE`$ a weak equivalence. A morphism $`(Y,h)(Z,g)`$ consists of a map of based spaces $`f:YZ`$ such that $`g\mathrm{\Sigma }^{\mathrm{}}f=h`$. Then define

$$𝔐_E:=|C_E|,$$

i.e., the geometric realization of (the nerve of) $`C_E`$.111There are set-theoretic difficulties presented by this definition, since $`C_E`$ isn’t a small category. The problem can be avoided in a variety of ways, all furnishing the same homotopy type (compare \[Wa, p. 379\]). We call $`𝔐_E`$ the moduli space of suspension structures on $`E`$.

Our third result gives a spectral sequence converging to the homotopy of $`𝔐_E`$ in positive degrees. Formulating it requires some preparation.

For $`q2`$, let $`W_q`$ be the spectrum with $`\mathrm{\Sigma }_q`$-action that classifies the $`q`$-th layer of the Goodwillie tower of the identity functor from based spaces to based spaces. Unequivariantly, $`W_q`$ is a wedge of $`(q1)`$ copies of the $`(1q)`$-sphere spectrum (see Johnson \[Jo\]). If $`q1`$, we take $`W_q`$ to be the trivial spectrum.

If $`E`$ is a spectrum, then its $`q`$-fold smash power $`E^q`$ has the structure of a spectrum with (naive) $`\mathrm{\Sigma }_q`$-action. Give the smash product $`W_qE^q`$ the diagonal $`\mathrm{\Sigma }_q`$-action. We are then entitled to form

$$W_q_{h\mathrm{\Sigma }_q}E^q,$$

i.e, the homotopy orbit spectrum of $`\mathrm{\Sigma }_q`$ acting on $`W_qE^q`$. Let $`F(E,W_q_{h\mathrm{\Sigma }_q}E^q)`$ denote the function space of spectrum maps from $`E`$ to $`W_q_{h\mathrm{\Sigma }_q}E^q`$. With these definitions we are ready to state our third result.

###### Theorem C (Spectral Sequence).

Let $`E`$ be $`1`$-connected. Assume $`𝔐_E`$ is non-empty and equipped with basepoint. Then there is a first quadrant spectral sequence converging to $`\pi _{p+1}(𝔐_E)`$ with

$$E_{p,q}^1:=\pi _p(F(E,W_q_{h\mathrm{\Sigma }_q}E^q))$$

and $`d^1`$-differential of bi-degree $`(1,1)`$.

Note when $`dimEn`$, we have $`E_{p,q}^1=0`$ for $`pq+1n`$. Thus for fixed $`p`$, $`E_{p,q}^1`$ is non-zero for finitely many $`q`$.

Our last result says that the obstruction $`\delta _E`$ is $`2`$-local.

###### Theorem D (Localization).

For any connected spectrum $`E`$, the homotopy class $`\delta _E`$ becomes trivial after inverting the prime $`2`$.

Furthermore, if $`f:EE^{}`$ is a map of spectra which is a $`2`$-local weak equivalence, then $`\delta _E=0`$ if and only if $`\delta _E^{}=0`$.

Consequently, if $`E`$ and $`E^{}`$ satisfy the connectivity assumptions of Theorem A, then $`E`$ is weak equivalent to a suspension spectrum if and only $`E^{}`$ is weak equivalent to a suspension spectrum.

### Remarks

Since the paper of Berstein and Hilton \[B-H\], the problem of deciding when a space has the homotopy type of a (single or iterated) suspension has been intensively studied.

To the best of my knowledge, the literature contains much less information about the desuspension problem for spectra, aside from the trivial stable range case (Freudenthal’s theorem) and a $`p`$-local version considered by Gray \[Gr\].

My original interest in the desuspension question for spectra came from embedding theory. In a future paper, we intend to use the above results to attack certain embedding questions.

###### Acknowledgements.

This paper would not have come into being had I not had the fortune of discussing mathematics with Greg Arone, Randy McCarthy and Tom Goodwillie in the middle to late 1990s. I am indebted to them for their help. I also owe much to Bill Richter, who explained to me the proof of the Berstein-Hilton theorem back in the 1980s.

I am grateful to Bob Bruner for explaining to me how to compute the stable homotopy of stunted projective space in low degrees, and for drawing my attention to the early work of Kuhn. I would like to thank Nick Kuhn for some very helpful remarks, especially concerning the proof of Lemma 7.1.

###### Outline.

§2 consists of preliminary material. The proof of Theorem A is contained in §3. In §4 we prove Theorem B. In §5 we express the connected components of the moduli space $`𝔐_E`$ as the classifying space of a suitable monoid. §6 contains the proof of Theorem C. Theorem D is proved in §7. In §8 we provide examples in connection with Theorem B. In §9 we discuss some loose ends. §10, which should be of independent interest, outlines the classification of quadratic homotopy functors from spectra to spectra. The classification presented here is due to Goodwillie, but I first learned about a special case of it from Kuhn. The classification gives a more concrete description of the invariant $`\delta _E`$.

## 2. Preliminaries

In this section we give the conventions and tools used throughout the paper.

### Spaces

All spaces will be compactly generated, and $`\mathrm{𝐓𝐨𝐩}`$ will denote the category of compactly generated spaces. In particular, we make the convention that products are to be retopologized with respect to the compactly generated topology. Let $`\mathrm{𝐓𝐨𝐩}_{}`$ denote the category of based spaces. A weak equivalence of spaces is shorthand for (a chain of) weak homotopy equivalence(s).

We use the usual connectivity terminology for spaces: A non-empty space is $`r`$-connected if its homotopy vanishes in degrees $`r`$ with respect to any choice of basepoint (note: every non-empty space is $`(1)`$-connected). A map $`AB`$ of spaces, with $`B`$ nonempty, is $`r`$-connected if for any choice of basepoint in $`B`$, the homotopy fiber with respect to this choice of basepoint is an $`(r1)`$-connected space.

A commutative square of spaces

(1)
$$\begin{array}{ccc}A& & C\\ & & & & \\ B& & D\end{array}$$

is said to be $`k`$-cocartesian if the map $`\text{hocolim}(BAC)D`$ is $`k`$-connected. It is $`\mathrm{}`$-cartesian if it is $`k`$-cartesian for all $`k`$. Similarly, the square is $`k`$-cartesian if the map $`A\text{holim }(BDC)`$ is $`k`$-connected. The square is $`\mathrm{}`$-(co)cartesian if it is $`k`$-(co)cartesian for every $`k`$.

### Spectra

A spectrum $`E`$ will be taken to mean a collection of based spaces $`\{E_i\}_i`$ together with based (structure) maps $`\mathrm{\Sigma }E_iE_{i+1}`$ where $`\mathrm{\Sigma }E_i`$ is the reduced suspension of $`E_i`$. A morphism of spectra $`EE^{}`$ consists of maps $`E_iE_i^{}`$ that are compatible with the structure maps. We denote the category of spectra by $`\mathrm{𝐒𝐩}`$.

A map of spectra is $`r`$-connected if it induces a surjection on homotopy up through degree $`r`$ and an isomorphism in degrees less than $`r`$. A spectrum is $`r`$-connected if the map to the zero object (consisting of the one point space in each degree) is $`(r+1)`$-connected.

A map of spectra is a weak equivalence if it is $`r`$-connected for all integers $`r`$. This notion of weak equivalence comes from a Quillen model structure on $`\mathrm{𝐒𝐩}`$. In this model structure, the fibrant objects are the $`\mathrm{\Omega }`$-spectra (those spectra $`E`$ such that $`E_i\mathrm{\Omega }E_{i+1}`$ is a weak equivalence for all $`i`$. A cofibrant object is (a retract of) a spectrum which is built up from the zero object by attaching cells. The model structure on $`\mathrm{𝐒𝐩}`$ comes equipped with functorial factorizations; in particular, fibrant and cofibrant approximation is functorial. For details, see Schwede \[Sc\].

We typically apply fibrant and/or cofibrant approximation functors to maintain homotopy invariance. To avoid clutter, we usually suppress the application of these approximations in the notation.

For example, if $`E`$ is a fibrant spectrum we write $`\mathrm{\Omega }^{\mathrm{}}E`$ for the zero space $`E_0`$ of $`E`$. If $`E`$ isn’t fibrant, to get a good construction we first replace $`E`$ by its associated $`\mathrm{\Omega }`$-spectrum $`E^{\mathrm{}}`$ and then define $`\mathrm{\Omega }^{\mathrm{}}E`$ as $`E_0^{\mathrm{}}`$.

If $`X`$ is a based space, its suspension spectrum $`\mathrm{\Sigma }^{\mathrm{}}X`$ has $`j`$-th space $`\mathrm{\Sigma }^jX`$, the $`j`$-fold reduced suspension of $`X`$ (for this to have the correct homotopy type, we assume that $`X`$ is a cofibrant space; i.e., the retract of a cell complex). In particular, $`Q(X):=\mathrm{\Omega }^{\mathrm{}}\mathrm{\Sigma }^{\mathrm{}}X`$ is the reduced stable homotopy functor. A map of spaces $`Y\mathrm{\Omega }^{\mathrm{}}E`$ is adjoint to a map of spectra $`\mathrm{\Sigma }^{\mathrm{}}YE`$.

The notion of $`r`$-cartesian and $`r`$-cocartesian squares generalizes to spectra in the obvious way. Note that an $`r`$-cartesian square is the same thing as and $`(r+1)`$-cocartesian square in the category of spectra. A commutative diagram such as (1) but this time of spectra, which is $`\mathrm{}`$-cartesian and with $`C`$ weak equivalent to the zero object, is called a fibration sequence (for spectra, a fibration sequence is the same thing as a cofibration sequence). By slight abuse of notation, we write the sequence as $`ABD`$.

As in the introduction, we write $`dimEn`$ if and only if $`E`$ is, up to homotopy, obtained from the zero object by attaching cells of dimension $`n`$. A spectrum is finite if it is built up from the zero object by attaching a finite number of cells. A spectrum is homotopy finite if it is weak equivalent to a finite spectrum.

In this paper, we can get away with a notion of smash product which is associative, commutative and unital up to homotopy (see e.g. Lewis, May and Steinberger \[L-M-S, Chap. 2\]). We also need to know that the extended power spectrum

$$D_k(E)=(E^k)_{h\mathrm{\Sigma }_k}:=(E\mathrm{\Sigma }_k)_+_{\mathrm{\Sigma }_k}E^k$$

exists, is functorial, homotopy invariant and coincides with the usual one in the case of spaces. For details, see \[L-M-S, Chap. 6\].

### Truncation

Let $`Y`$ be a based space, $`W`$ be a $`1`$-connected spectrum and $`f:\mathrm{\Sigma }^{\mathrm{}}YW`$ a map of spectra. Assume $`dimWn`$ and that $`f`$ is $`n`$-connected.

###### Lemma 2.1.

If $`n0`$, there exists a space $`Z`$, and an $`(n1)`$-connected map $`g:ZY`$ such that the composite

$$\mathrm{\Sigma }^{\mathrm{}}Z\stackrel{\mathrm{\Sigma }^{\mathrm{}}g}{}\mathrm{\Sigma }^{\mathrm{}}Y\stackrel{𝑓}{}W$$

is a weak equivalence.

###### Proof.

If $`n1`$, then $`W`$ is weak equivalent to the zero object, and the proof is trivial in this case.

Assume then that $`n2`$. Let $`h`$ be the Eilenberg-Mac Lane spectrum. Then the induced map

$$h\mathrm{\Sigma }^{\mathrm{}}YhW$$

can be thought of as a map of chain complexes (via the Dold-Kan correspondence). Applying the truncation lemma \[Kl, 4.1\], We obtain a space $`Z`$ and an $`(n1)`$-connected map $`ZY`$ such that the composite

$$h\mathrm{\Sigma }^{\mathrm{}}Zh\mathrm{\Sigma }^{\mathrm{}}YhW$$

is a weak equivalence. But this composite is obtained by smashing the composite

$$\mathrm{\Sigma }^{\mathrm{}}Z\mathrm{\Sigma }^{\mathrm{}}YW$$

with $`h`$. The result then follows by Whitehead’s theorem. ∎

###### Corollary 2.2 (Freudenthal).

Let $`E`$ be a an $`r`$-connected spectrum, $`r1`$. Assume $`dimE2r+2`$. Then $`E`$ is weak equivalent to a suspension spectrum.

###### Proof.

This follows from 2.1 because the map $`\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}EE`$ is $`(2r+2)`$-connected (see 3.2 below). ∎

### Another description of the moduli space

One can use Lemma 2.1 to give an alternative description of the moduli space when $`E`$ is $`1`$-connected. Let $`C_E^{}`$ be the (not full) subcategory of $`C_E`$ having the same objects, and whose morphisms $`(Y,h)(Y^{},h^{})`$ satisfy the condition that $`YY^{}`$ is a weak homotopy equivalence. Let $`𝔐_E^{}`$ denote the realization of $`C_E^{}`$.

###### Proposition 2.3.

Assume that $`E`$ is fibrant, cofibrant and $`1`$-connected. Then the inclusion $`𝔐_E^{}𝔐_E`$ is a homotopy equivalence.

###### Proof.

Let $`YY^+`$ denote the plus construction. Applying functorial factorization, we can arrange it so that the natural map $`YY^+`$ is a cofibration.

If $`(Y,h)`$ is an object of $`C_E`$, then the $`1`$-connectedness of $`E`$ implies $`Y`$ has trivial first homology. Therefore, $`Y^+`$ is a $`1`$-connected space. So if $`(Y,h)(Z,h^{})`$ is a morphism, the induced map of plus constructions $`Y^+Z^+`$ is a weak equivalence by Whitehead’s theorem.

Since $`E`$ is fibrant and cofibrant, the space $`\mathrm{\Omega }^{\mathrm{}}E`$ is cofibrant and $`1`$-connected. The natural map $`\mathrm{\Omega }^{\mathrm{}}E(\mathrm{\Omega }^{\mathrm{}}E)^+`$ is a cofibration which is a weak equivalence. Consequently, there is a retraction $`r:(\mathrm{\Omega }^{\mathrm{}}E)^+\mathrm{\Omega }^{\mathrm{}}E`$.

Let $`i:C_E^{}C_E`$ denote the inclusion. Define a functor

$$\pi :C_EC_E^{}$$

by the rule $`\pi (Y,h)=(Y^+,h^{\mathrm{}})`$, where $`h^{\mathrm{}}:\mathrm{\Sigma }^{\mathrm{}}(Y^+)E`$ is adjoint to the composite

$$Y^+\stackrel{(\widehat{h})^+}{}(\mathrm{\Omega }^{\mathrm{}}E)^+\stackrel{𝑟}{}\mathrm{\Omega }^{\mathrm{}}E.$$

Then there is an evident natural transformation from the identity functor of $`C_E`$ to the composite functor $`i\pi `$. There is a similar evident natural transformation from the identity functor of $`C_E^{}`$ to $`\pi i`$. It follows that $`i`$ induces a homotopy equivalence on realizations. ∎

## 3. Proof of Theorem A

After defining the obstruction $`\delta _E`$, the idea of the remainder of the proof will be to construct a highly connected map from a suspension spectrum to $`E`$. The proof is then completed by applying Lemma 2.1.

### Definition of the obstruction

We describe below a certain fibration sequence of spectra

$$D_2(E)S_2(E)E$$

in which $`S_2:\mathrm{𝐒𝐩}\mathrm{𝐒𝐩}`$ is a certain homotopy functor.

Assuming this construction has been specified, we have:

###### Definition 3.1.

The class

$$\delta _E[E,\mathrm{\Sigma }D_2(E)]$$

is the obstruction to splitting the above fibration sequence, i.e., the homotopy class of its connecting map to the right.

### Construction of the fibration

Consider the homotopy functor

$$\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}:\mathrm{𝐒𝐩}\mathrm{𝐒𝐩}$$

which assigns to a spectrum the suspension spectrum of its zero space. Let $`S_k(E)`$ denote the $`k`$-th stage of the Goodwillie tower of this functor and let $`F_k(E):=\text{fiber}(S_k(E)S_{k1}(E))`$ denote the $`k`$-th layer.

The following result has been noted by several people, including Goodwillie, Arone, McCarthy, and Ahearn and Kuhn \[A-K, Cor. 1.3\].

###### Lemma 3.2.

Assume $`E`$ is $`r`$-connected. Then the map

$$\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}ES_k(E)$$

is $`((k+1)r+k+1)`$-connected. Consequently, the Goodwillie tower of $`\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}E`$ is convergent if $`E`$ is $`0`$-connected. Furthermore, there is a natural weak equivalence of functors

$$F_k(E)D_k(E).$$

Applying 3.2, we see that the bottom of the tower yields a fibration sequence of spectra

$$D_2(E)S_2(E)E$$

together with a $`(3r+3)`$-connected map $`\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}ES_2(E)`$.

###### Corollary 3.3.

Assume $`E`$ is $`r`$-connected and $`dimE3r+3`$. Then $`\delta _E=0`$ if and only if the map $`\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}EE`$ admits a section up to homotopy.

###### Proof.

The class $`\delta _E`$ is trivial if and only if the $`S_2(E)E`$ has a section. The dimension constraint on $`E`$ and the lemma show that $`S_2(E)E`$ admits a section if and only if $`\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}EE`$ admits a section up to homotopy. ∎

###

Now assume $`E`$ is $`r`$-connected, $`r1`$, $`dimE3r+3`$ and $`\delta _E=0`$. By the above remarks, we are entitled to choose a homotopy section $`\sigma :E\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}E`$. Let

$$\mathrm{\Omega }^{\mathrm{}}\sigma :\mathrm{\Omega }^{\mathrm{}}E\mathrm{\Omega }^{\mathrm{}}\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}E$$

be the corresponding map of zero spaces.

There is another map $`c:\mathrm{\Omega }^{\mathrm{}}E\mathrm{\Omega }^{\mathrm{}}\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}E`$ (not an infinite loop map) which is defined by taking the adjoint to the identity map of $`\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}E`$.

Let $`Y`$ be the homotopy pullback of the maps $`\mathrm{\Omega }^{\mathrm{}}\sigma `$ and $`c`$. Thus we have an $`\mathrm{}`$-cartesian square of spaces

(2)
(commutative up to preferred homotopy).

###### Lemma 3.4.

With respect to the above assumptions, let $`\widehat{j}:\mathrm{\Sigma }^{\mathrm{}}YE`$ be the adjoint to the map labeled $`j`$ in diagram (2). Then $`\widehat{j}`$ is $`(3r+2)`$-connected.

###### Proof.

Consider the diagram of spectra

(3)
in which the maps labeled with $`q`$ are the counits to the adjunctions for $`E`$ and $`\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}E`$. The map $`\widehat{j}`$ is therefore given by the composite $`q_E\mathrm{\Sigma }^{\mathrm{}}j`$, and the composite along the bottom, $`q_{\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}E}\mathrm{\Sigma }^{\mathrm{}}c`$, is clearly the identity. The left square commutes up to a preferred homotopy. The right square commutes on the nose.

The maps $`\mathrm{\Omega }^{\mathrm{}}\sigma `$ and $`c`$ in diagram (2) are both $`(2r+1)`$-connected. Consequently, by the dual Blakers-Massey theorem (see \[Go2\]) diagram (2) is $`(4r+3)`$-cocartesian. In particular, the left square in diagram (3) is a $`(4r+3)`$-cocartesian square of spectra.

Regarding the right square in diagram (3), If $`C`$ denotes the homotopy cofiber of $`\sigma `$, then the evident map $`CD_2(E)`$ is $`(3r+3)`$-connected. Similarly, if $`C^{}`$ denotes the homotopy cofiber of $`\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}\sigma `$, an argument using the Blakers-Massey theorem, which we omit, shows that the evident map $`C^{}\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}D_2(E)`$ is also $`(3r+3)`$-connected.

Corollary 2.2 shows that the map $`\mathrm{\Sigma }^{\mathrm{}}\mathrm{\Omega }^{\mathrm{}}D_2(E)D_2(E)`$ is $`(4r+4)`$-connected (because $`D_2(E)`$ is $`(2r+1)`$-connected). We infer that the map $`C^{}C`$ is also $`(3r+3)`$-connected. Therefore the right square in diagram (3) is $`(3r+3)`$-cocartesian.

Putting both squares together, it follows that diagram (3) is $`(3r+3)`$-cocartesian. Since the bottom composite is the identity map, we infer that the top composite $`\widehat{j}`$ is $`(3r+2)`$-connected, as asserted. ∎

### Completion of the proof

Assume in addition to the above that $`dimE3r+2`$. Since the map $`\widehat{j}:\mathrm{\Sigma }^{\mathrm{}}YE`$ is $`(3r+2)`$-connected we can apply 2.1. This gives a based space $`Z`$ and a based map $`ZY`$ such that the composite $`\mathrm{\Sigma }^{\mathrm{}}Z\mathrm{\Sigma }^{\mathrm{}}YE`$ is a weak equivalence. The proof of Theorem A is now complete.

## 4. Proof of Theorem B

### Step 1

Let $`𝒮_E`$ denote the set of homotopy classes of sections of the fibration $`S_2(E)E`$. Note that the abelian group $`[E,D_2(E)]`$ acts freely and transitively on $`𝒮_E`$ (cf. Lemma 6.1 below). Thus if $`𝒮_E`$ is given a basepoint, it follows that $`[E,D_2(E)]`$ and $`𝒮_E`$ are isomorphic (the isomorphism is dependent on the choice of basepoint).

There is a function

$$\varphi :\mathrm{\Theta }_E𝒮_E$$

defined by sending $`(Y,h)`$ to the class represented by

$$E\mathrm{\Sigma }^{\mathrm{}}YS_2(\mathrm{\Sigma }^{\mathrm{}}Y)S_2(E)$$

(the map in the middle is the preferred section of $`S_2(\mathrm{\Sigma }^{\mathrm{}}Y)\mathrm{\Sigma }^{\mathrm{}}Y`$).

If $`ES_2(E)`$ is a representative of $`𝒮_E`$, and $`dimE3r+2`$, we can perform the constructions in the previous section to get an element of $`\mathrm{\Theta }_E`$ (the element is not necessarily unique). It is straightforward to check that this element of $`\mathrm{\Theta }_E`$ maps to the given element of $`𝒮_E`$. Thus, when $`dimE3r+2`$, the function $`\varphi :\mathrm{\Theta }_E𝒮_E`$ is onto.

To complete the proof of Theorem B, it will be sufficient to show that $`\varphi `$ is one-to-one when $`dimE3r+1`$. Choose a basepoint $`(Y,h)`$ for $`\mathrm{\Theta }_E`$. Then $`𝒮_E`$ inherits a basepoint and $`𝒮_E`$ becomes identified with $`[E,D_2(E)]`$. Thus we may rewrite $`\varphi `$ as a basepoint preserving function

$$\mathrm{\Theta }_E[E,D_2(E)],$$

where the basepoint of the codomain is the zero element. Since $`Y`$ was chosen arbitrarily, it suffices to show that $`\varphi `$ is one-to-one at the inverse image of the basepoint.

### Step 2

We digress to develop a relative version of Theorem A. Suppose that $`AE`$ is a cofibration in the category of spectra. We write $`dim(E,A)n`$ if $`E`$ is obtained from $`A`$ up to homotopy by attaching cells of dimension $`n`$. Assume that $`A=\mathrm{\Sigma }^{\mathrm{}}Z`$ is a suspension spectrum.

###### Theorem 4.1.

There is an obstruction

$$\delta _{(E,A)}[E/A,\mathrm{\Sigma }D_2(E)]$$

whose triviality is necessary to finding a cofibration $`ZW`$ and a weak equivalence $`\mathrm{\Sigma }^{\mathrm{}}WE`$ extending the identity on $`\mathrm{\Sigma }^{\mathrm{}}Z=A`$. If $`E`$ is $`r`$-connected ($`r1`$) and $`dim(E,A)3r+2`$, then the vanishing of this obstruction is also sufficient.

The obstruction $`\delta _{(E,A)}`$ is defined in the same way as $`\delta _E`$ taking care to notice that the restriction of $`\delta _E`$ to $`A`$ has a preferred trivialization. The proof of 4.1 is virtually the same as the proof of Theorem A, so we omit it.

### Step 3

We now return to the proof of Theorem B. We shall apply 4.1 in the following situation: choose a representative $`(Y,h)`$ for an element in $`\mathrm{\Theta }_E`$ (call this the basepoint). Suppose that $`(Y^{},h^{})`$ represents an another element in $`\mathrm{\Theta }_E`$. Both elements combine to give a weak equivalence

$$\mathrm{\Sigma }^{\mathrm{}}(YY^{})\stackrel{}{}EE.$$

Applying 4.1 to the pair $`(EI_+,EE)`$ we get a necessary obstruction

$$\delta _{(EI_+,EE)}[\mathrm{\Sigma }E,\mathrm{\Sigma }D_2(E)]=[E,D_2(E)]$$

to finding a cofibration $`YY^{}W`$ of based spaces and a weak equivalence

$$(\mathrm{\Sigma }^{\mathrm{}}W,\mathrm{\Sigma }^{\mathrm{}}(YY^{}))(EI_+,EE).$$

It follows from this that the zig-zag $`YWY^{}`$ equates $`(Y,h)`$ with $`(Y^{},h^{})`$ in $`\mathrm{\Theta }_E`$.

Thus, fixing $`(Y,h)`$ and allowing $`(Y^{},h^{})`$ to vary defines a function

$$\mathrm{\Theta }_E[E,D_2(E)],$$

which is just another description of the function $`\varphi `$ (we omit the details).

Assume that $`E`$ is $`r`$-connected and $`dimE3r+1`$. Then 4.1 shows that $`(Y^{},h^{})`$ maps to zero under $`\varphi `$ if and only if $`(Y^{},h^{})`$ is the basepoint of $`\mathrm{\Theta }_E`$. This completes the proof of Theorem B.

## 5. $`𝔐_E`$ as a classifying space

Fix a cofibrant based space $`Z`$. Let $`\mathrm{𝐓𝐨𝐩}_{/Z}`$ denote the category of based spaces over $`Z`$. An object $`y`$ of this category consists of a based space $`Y`$ together with a based map $`p_Y:YZ`$. A morphism $`(Y,p_Y)(Y^{},p_Y^{})`$ is given by a map of based spaces $`f:YY^{}`$ such that $`p_Y^{}f=p_Y`$.

Since $`\mathrm{𝐓𝐨𝐩}_{/Z}`$ is an over category of the Quillen model category $`\mathrm{𝐓𝐨𝐩}_{}`$, it follows that it too has the structure of a Quillen model category (see \[Qu\]). A weak equivalence is a morphism $`yy^{}`$ whose underlying map of spaces is a weak homotopy equivalence. We say $`yy^{}`$ is a fibration if its underlying map of spaces is. We say that $`yy^{}`$ is a cofibration if it satisfies the left lifting property with respect to the acyclic fibrations.

Let $`w\mathrm{𝐓𝐨𝐩}_{/Z}`$ denote the subcategory consisting of the weak equivalences. Let $`w_{(y)}\mathrm{𝐓𝐨𝐩}_{/Z}`$ denote the full subcategory of $`w\mathrm{𝐓𝐨𝐩}_{/Z}`$ consisting of those objects connected to $`y`$ by a chain of weak equivalences.

The proof of the following proposition, which we attribute to Waldhausen, is proved by the same method as \[Wa, 2.2.5\]. We omit the details.

###### Proposition 5.1.

Assume that $`y`$ is fibrant and cofibrant. Then there is a weak equivalence of spaces

$$|w_{(y)}\mathrm{𝐓𝐨𝐩}_{/Z}|BG(y),$$

where $`G(y)`$ denotes the topological monoid of self homotopy equivalences of $`y`$ in $`\mathrm{𝐓𝐨𝐩}_{/Z}`$, and $`BG(y)`$ is its classifying space.

We apply this in the following special case: let $`E`$ be a $`1`$-connected fibrant and cofibrant spectrum. The category $`C_E^{}`$, whose realization is the moduli space $`𝔐_E^{}`$ (cf. §2), is just the full subcategory of $`w\mathrm{𝐓𝐨𝐩}_{/\mathrm{\Omega }^{\mathrm{}}E}`$ whose objects $`Y`$ are such that the map $`p_Y:Y\mathrm{\Omega }^{\mathrm{}}E`$ is adjoint to a weak equivalence of spectra.

If we combine Proposition 5.1 with Proposition 2.3, we obtain

###### Corollary 5.2.

Let $`y`$ be an object of $`C_E`$ which is fibrant and cofibrant when considered as an object of $`\mathrm{𝐓𝐨𝐩}_{/\mathrm{\Omega }^{\mathrm{}}E}`$. Let $`𝔐_{E,(y)}`$ be the connected component of $`𝔐_E`$ which contains $`y`$. Then there is a homotopy equivalence

$$𝔐_{E,(y)}BG(y).$$

## 6. Proof of Theorem C

### Outline of the proof

The basepoint $`y=(Y,h)`$ of $`𝔐_E`$ can be taken as an object of $`C_E`$. Taking the plus construction if necessary, we can assume without loss in generality that $`Y`$ is $`1`$-connected (see the argument in the proof of Proposition 2.3). We may also assume that $`y`$ is fibrant and cofibrant when considered as an object of $`\mathrm{𝐓𝐨𝐩}_{/\mathrm{\Omega }^{\mathrm{}}E}`$.

We will construct a tower of fibrations of based spaces

$$\mathrm{}T_3(y)T_2(y)T_1(y)$$

such that

* $`T_1(y)`$ is contractible;
* for $`k>1`$, there is a weak equivalence

$$\text{fiber}(T_k(y)T_{k1}(y))\mathrm{\Omega }^{\mathrm{}}(W_k_{h\mathrm{\Sigma }_k}Y^k),$$

where $`W_k`$ denotes the spectrum that classifies the $`k`$-th layer of the Goodwillie tower of the identity functor on based spaces.
* there is a weak equivalence

$$\mathrm{\Omega }𝔐_{E,(y)}\underset{k}{lim}T_y(y).$$

Assuming this has been done, we can define the spectral sequence $`\{E_{p,q}^r\}`$ as the homotopy spectral sequence of the tower $`\{T_k(y)\}_k`$.

We now digress to discuss generalities about section spaces and the basic properties of the Goodwillie tower of the identity functor.

### Digression

Suppose $`p:EZ`$ is a fibration of based spaces. We say that $`p`$ is principal if there exists a commutative $`\mathrm{}`$-cartesian square of based spaces

$$\begin{array}{ccc}E& & P\\ p& & & & \\ Z& & B\end{array}$$

such that $`P`$ is contractible.

Suppose that $`Z`$ is connected. If $`p:EZ`$ is principal, there is an “action” $`\mathrm{\Omega }B\times EE`$. If there is a section $`ZE`$ one can combine it with this action to produce a map of fibrations $`\mathrm{\Omega }B\times ZE`$ covering the identity map of $`Z`$. This implies that $`p`$ is weak fiber homotopically trivial. Let $`\text{Sec}(p)`$ denote the space of sections of $`p`$. Then we have shown

###### Lemma 6.1.

Assume $`p:EZ`$ is principal. Assume that $`\text{Sec}(p)`$ is non-empty and comes equipped with basepoint. Then there is a weak equivalence of based spaces

$$\text{Sec}(p)F(Z,\mathrm{\Omega }B).$$

We next recall for the reader the basic properties of the Goodwillie tower of the identity functor on based spaces (cf. Goodwillie \[Go1\], \[Go2\], \[Go3\], Johnson \[Jo\]).

###### Theorem 6.2.

There is a tower of fibrations of homotopy functors on based spaces

$$\mathrm{}P_2(X)P_1(X)$$

and compatible natural transformations $`XP_k(X)`$ such that

* $`P_1(X)=Q(X)`$ is the stable homotopy functor;
* For $`k2`$, the fibration $`P_k(X)P_{k1}(X)`$ is principal;
* the $`k`$-the layer $`L_k(X)`$ is naturally weak equivalent to the functor

$$X\mathrm{\Omega }^{\mathrm{}}(W_k_{h\mathrm{\Sigma }_k}X^k);$$
* if $`X`$ is $`1`$-connected, then the natural map

$$X\underset{k}{lim}P_k(X)$$

is a weak equivalence.

### Completion of the proof

We are now in a position to define the tower $`\{T_k(y)\}_k`$.

###### Definition 6.3.

Let $`T_k(y)`$ be the space of lifts

where $`YP_1(Y)`$ is the natural map. Note that $`T_k(y)`$ comes equipped with a basepoint defined by the natural map $`YP_k(Y)`$.

From the definition of $`T_k(y)`$, there is an evident tower of fibrations of based spaces

$$\mathrm{}T_2(y)T_1(y)$$

with $`T_1(y)=`$. Furthermore, the fiber of the map $`T_k(y)T_{k1}(y)`$ at the basepoint is identified with the space of lifts

Note that this is just the space of sections of the pulled back fibration $`Y\times _{P_{k1}(Y)}P_k(Y)Y`$. The latter fibration is principal and comes equipped with a preferred section.

Applying 6.1 and 6.2 we obtain a weak equivalence of based spaces

$$\text{fiber}(T_k(y)T_{k1}(y))F(Y,\mathrm{\Omega }^{\mathrm{}}(W_k_{h\mathrm{\Sigma }_k}Y^k)).$$

Taking adjunctions, we see that $`F(Y,\mathrm{\Omega }^{\mathrm{}}(W_k_{h\mathrm{\Sigma }_k}Y^k))`$ is weak equivalent to the space $`F(E,W_k_{h\mathrm{\Sigma }_k}E^k)`$ (where we are using the identification $`\mathrm{\Sigma }^{\mathrm{}}YE`$).

To complete the proof of Theorem C, it suffices to identify the inverse limit of the tower $`\{T_k(y)\}_k`$. It is clear that this inverse limit is just the space of lifts

Recall that $`P_1(Y)=Q(Y)`$ and since $`Y`$ is $`1`$-connected, $`lim_kP_k(Y)Y`$. With respect to these identifications, we infer that this space of lifts is weak equivalent to the underlying space of the topological monoid $`G(y)`$. Applying 5.2 and 2.3 completes the proof of Theorem C.

## 7. Proof of Theorem D

Let $`\text{tr}:D_2(E)E^2`$ be the transfer and let $`\pi :E^2D_2(E)`$ be the projection (\[L-M-S, Chap. 4\]). It is well-known that the composite

$$p\text{tr}:D_2(E)D_2(E)$$

is a weak equivalence after inverting $`2`$. To prove the first part of Theorem D, it is clearly enough to show

###### Lemma 7.1.

For a connected spectrum $`E`$, the class

$$(\mathrm{\Sigma }\text{tr})\delta _E[E,\mathrm{\Sigma }EE]$$

is trivial.

###### Proof.

This is a consequence a result of Ahearn and Kuhn (see \[A-K, Cor. 1.7\] in conjunction with the discussion after \[A-K, Cor. 1.2\]). For an alternative proof, see the remark after Corollary 10.5 below.

In the case we need, Ahearn and Kuhn construct a natural transformation $`\mathrm{\Phi }:S_2(E)EE`$ fitting into a preferred homotopy commutative diagram

They then identify the composite map

$$D_2(E)S_2(E)\stackrel{\Phi }{}EE$$

with the transfer.222Randy McCarthy has informed me that that this result is also implicit in the recent Ph.D. thesis of K.B. Bauer \[Ba\].

Since $`EE^2`$ is a homogeneous quadratic functor, it coincides with the second stage of its Goodwillie tower and its first stage is trivial. So the map of Goodwillie towers associated with $`\mathrm{\Phi }`$ in degrees $`2`$ may be displayed as

If we take horizontal homotopy cofibers of the maps from stage two to stage one, we obtain a homotopy commutative diagram

which shows that $`(\mathrm{\Sigma }\mathrm{tr})\delta _E`$ is null homotopic. ∎

We now prove the second part of Theorem D. Let $`f:EE^{}`$ be a map which is a $`2`$-local weak equivalence. Assume that $`E`$ and $`E^{}`$ satisfy the connectivity hypotheses of Theorem B. Then the naturality of the obstruction shows that the diagram

$$\begin{array}{ccc}E& \stackrel{\delta _E}{}& \mathrm{\Sigma }D_2(E)\\ f& & \mathrm{\Sigma }D_2\left(f\right)& & \\ E^{}& \underset{\delta _E^{}}{}& \mathrm{\Sigma }D_2(E^{})\end{array}$$

is homotopy commutative. If we invert $`2`$, then by the first part, the obstructions $`\delta _E`$ and $`\delta _E^{}`$ vanish. It we localize at 2, then $`f`$ and $`\mathrm{\Sigma }D_2(f)`$ become weak equivalences, and therefore $`\delta _E`$ and $`\delta _E^{}`$ coincide. Thus, integrally, $`\delta _E=0`$ if and only $`\delta _E^{}=0`$.

## 8. Illustrations of Theorem B

We give two examples.

### Two cells

For $`p>1`$ and $`q3p3`$, consider the suspension spectrum $`E=\mathrm{\Sigma }^{\mathrm{}}(S^p_fe^{q+1})`$ where $`f:S^qS^p`$ is some map. According to Theorem B, we have

$$\mathrm{\Theta }_E[E,D_2(E)]\pi _{q+1}(D_2(S^p)),$$

where the second of these isomorphisms comes from that fact that a homotopy class $`ED_2(E)`$ maps the top cell of $`E`$ into $`D_2(S^p)D_2(E)`$.

The first non-trivial group occurs when $`q=2p1`$. The group $`\pi _{2p}(D_2(S^p))`$ is isomorphic to $``$ if $`p`$ is even and $`_2`$ if $`p`$ is odd. The distinct elements of $`\mathrm{\Theta }_E`$ are represented by the suspension spectra of $`S^p_ge^{2p}`$, with

$$g=f+k[\iota ,\iota ],$$

where $`k`$ is an integer (mod 2 if $`p`$ is odd), and $`[\iota ,\iota ]`$ is the Whitehead product. The group structure on $`\mathrm{\Theta }_E`$ is given by adding these integers.

In fact, at the prime 2, Mahowald has computed $`\pi _{q+1}(D_2(S^p))`$ for $`q\mathrm{min}(3p3,2p+29)`$ (see \[Mah, table 4.1\], see also Milgram \[Mi, table 13.5\]). For example, assuming in addition $`p1\text{ mod }16`$, the first few groups are

$$\begin{array}{ccccccccc}& j& 0& \text{1}& \text{2}& \text{3}& \text{4}& \text{5}& \text{6}\\ & \pi _{2p+j}(D_2(S^p))& _2& _2& _8& _2& 0& _2& _{16}_2\end{array}.$$

### Tori

Let $`E=\mathrm{\Sigma }^{\mathrm{}}(S^p\times S^p)`$, with $`p>1`$. Then

$$\mathrm{\Theta }_E[E,D_2(E)]\pi _{2p}(D_2(S^pS^p)).$$

When $`p`$ is even, the group $`\pi _{2p}(D_2(S^pS^p))`$ is isomorphic to $`^3`$. When $`p`$ is odd, it is isomorphic to $`_2_2`$. The elements of $`\mathrm{\Theta }_E`$ can be represented by the suspension spectra of the complexes

$$(S^pS^p)_he^{2p}$$

with attaching map

$$h=kx[\iota ,\iota ]+\mathrm{}y[\iota ,\iota ]+m\omega ,$$

where $`\omega `$ is the attaching map for the top cell of $`S^p\times S^p`$, $`x`$ and $`y`$ are the summand inclusions $`S^pS^pS^p`$, and $`k,\mathrm{}`$ and $`m`$ are integers (take $`k`$ and $`\mathrm{}`$ modulo 2 if $`p`$ is odd). Hence the elements of $`\mathrm{\Theta }_E`$ are specified by triples $`(k,\mathrm{},m)`$. The identity element of $`\mathrm{\Theta }_E`$ is $`(0,0,1)`$ and addition is given by

$$(k,\mathrm{},m)+(k^{},\mathrm{}^{},m^{})=(k+k^{},\mathrm{}+\mathrm{}^{},m+m^{}1).$$

## 9. Loose Ends

### Relation to the James-Hopf invariant

Another approach to desuspension questions is to inductively desuspend cell-by-cell (with respect to a cell decomposition of $`E`$). We explain here how this relates to our approach.

The idea is this: when $`E`$ can be written as a homotopy cofiber of a map of suspension spectra $`f:\mathrm{\Sigma }^{\mathrm{}}A\mathrm{\Sigma }^{\mathrm{}}B`$, it turns out that $`\delta _E`$ is closely related to the James-Hopf invariant

$$H_2(f)[\mathrm{\Sigma }^{\mathrm{}}A,D_2(\mathrm{\Sigma }^{\mathrm{}}B)].$$

Explicitly, there is a homomorphism

$$\psi :[\mathrm{\Sigma }^{\mathrm{}}A,D_2(\mathrm{\Sigma }^{\mathrm{}}B)][E,\mathrm{\Sigma }D_2(E)]$$

given by suspending, precomposing with the connecting map $`E\mathrm{\Sigma }\mathrm{\Sigma }^{\mathrm{}}A`$ and postcomposing with the inclusion $`\mathrm{\Sigma }D_2(\mathrm{\Sigma }^{\mathrm{}}B)\mathrm{\Sigma }D_2(E)`$. We assert that $`\psi `$ maps $`H_2(f)`$ to $`\delta _E`$ (we defer the proof to another paper).

If we also assume that $`B`$ is $`r`$-connected ($`r1`$), $`dimB2r+1`$, $`A`$ is $`2r`$-connected and $`dimA3r+1`$, then we infer that $`E`$ is $`r`$-connected and $`dimE3r+2`$. Furthermore, obstruction theory implies that $`\psi `$ is an isomorphism. We conclude that $`E`$ desuspends if and only if $`H_2(f)=0`$ (the ‘if’ part of this statement is well-known).

### Musings on the spectral sequence

It would be desirable to have a version of the spectral sequence in Theorem C which converges to $`\pi _{}(𝔐_E)`$ in degree zero.

The spectral sequence still lacks a geometric interpretation. Conjecturally, the spectral sequence should be a packaging machine for the obstructions to equipping a spectrum with the structure of an

$`E_{\mathrm{}}`$-coalgebra over the sphere spectrum.”

What I have in mind here comes from an (unproved) observation that $`S_2(E)`$ is a model for $`(EE)^_2=`$ the categorical fixed points of $`_2`$ acting on $`EE`$, and a choice of section for the fibration $`S_2(E)E`$ can be thought of as commutative (but not necessarily associative) “diagonal” for $`E`$.

It is tempting to conjecture that the spectral sequence arises from a tower whose $`k`$-stage encodes, in some sense, the moduli space of spectra equipped with choice of commutative diagonal that is coherently homotopy associative up to order $`k1`$.

## 10. Appendix: On quadratic homotopy functors

The invariant $`\delta _E`$ was defined as the connecting map associated with the degree two homotopy functor $`ES_2(E)`$. It is not difficult to classify the quadratic homotopy functors from spectra to spectra. The classification, which is expressed in terms of the Tate construction, gives a more concrete way of understanding $`\delta _E`$. I would like to thank Tom Goodwillie and Nick Kuhn for explaining these ideas to me.

Let $`f:\mathrm{𝐒𝐩}\mathrm{𝐒𝐩}`$ be a homotopy functor of degree $`2`$. This means that $`f`$ maps strongly cocartesian $`3`$-cubes to $`\mathrm{}`$-cartesian (= $`\mathrm{}`$-cocartesian) $`3`$-cubes. We also require $`f`$ to satisfy the usual technical conditions appearing in the functor calculus. Namely, we assume that $`f`$ preserves filtered homotopy colimits (in particular, $`f`$ is determined up to weak equivalence by its restriction to finite spectra) and is stably excisive (see \[Go2\] for the meaning of this language). We also assume that $`f`$ is reduced in the sense that its value at the zero object is contractible.

The general theory describes $`f`$ up to homotopy as sitting within a fibration sequence

$$W_2_{h_2}E^2f(E)W_1E,$$

where $`W_1`$ is a spectrum and $`W_2`$ is a spectrum with $`_2`$-action (we remind the reader that $`W_2_{h_2}E^2`$ denotes the homotopy orbits of $`_2`$ acting diagonally on $`W_2E^2`$). The connecting map of this sequence is then a natural transformation

$$u_f(E):W_1E\mathrm{\Sigma }(W_2_{h_2}E^2).$$

Renaming $`W:=W_1`$ and $`V:=\mathrm{\Sigma }W_2`$. We can rewrite $`u_f(E)`$ as

$$WEV_{h_2}E^2.$$

Notice that $`u_f(E)`$ is a natural transformation from the linear functor $`EWE`$ to the homogeneous quadratic functor $`EV_{h_2}E^2`$, and clearly, $`f(E)`$ is identified with the homotopy fiber of $`u_f(E)`$. We will show how to classify such natural transformations.

Let $`\alpha `$ denote the unique non-trivial one dimensional real representation of $`_2`$. If $`j0`$ is any integer, let $`S^{j\alpha }`$ be the suspension spectrum of the one point compactification of the direct sum of $`j`$-copies of $`\alpha `$. If $`j<0`$, define $`S^{j\alpha }`$ as $`F(S^{j\alpha },S^0)`$, the Spanier-Whitehead dual of $`S^{j\alpha }`$. Note that $`S^{j\alpha }`$ has a $`_2`$-action, and there is an evident cofibration sequence of spectra with $`_2`$-action

$$S^{(j1)\alpha }S^{j\alpha }(_2)_+S^j$$

where the third term is given by inducing up the $`j`$-sphere with trivial action to a spectrum with free $`_2`$-action.

Let us now consider $`u_f`$ applied to $`S^jE`$. Then we have

$$u_f(S^jE):WES^j(VE^2)_{h_2}(S^jS^{j\alpha })$$

where we have taken the liberty of rewriting $`S^jS^j`$ with tranposition action as $`S^jS^{j\alpha }`$. If we desuspend $`j`$-times, we get for each $`j`$ a map

$$\mathrm{\Sigma }^ju_f(S^jE):WE(VE^2)_{h_2}S^{j\alpha }.$$

###### Lemma 10.1.

These maps are compatible in the sense that the diagram

is homotopy commutative.

###### Proof.

If we apply $`u_f`$ to the evident cocartesian square of spectra

$$\begin{array}{ccc}S^jE& \stackrel{i_+}{}& D^{j+1}E\\ i_{}& & & & \\ D^{j+1}E& & S^{j+1}E\end{array}$$

we obtain a map of squares of spectra. This map induces a homotopy commutative diagram

$$\begin{array}{ccc}\mathrm{\Sigma }(WES^j)& \stackrel{\mathrm{\Sigma }u_f(S^jE)}{}& \mathrm{\Sigma }(VE^2_{h_2}(S^jS^{j\alpha }))\\ & & & & \\ WES^{j+1}& \underset{u_f(S^{j+1}E)}{}& VE^2_{h_2}(S^{j+1}S^{(j+1)\alpha })\end{array}$$

where the vertical maps are induced by maps from the pushouts coming from $`i_\pm `$. The proof is now completed by desuspending $`j+1`$ times. ∎

Recall that any spectrum $`U`$ with (naive) $`_2`$-action comes equipped with a norm map $`U_{h_2}U^{h_2}`$ from homotopy orbits to homotopy fixed points (for details, see \[A-C-D\], \[W-W, §2\], \[G-M\], \[Kl\]). Recall that the Tate spectrum

$$U^{t_2}$$

is the homotopy cofiber of the norm map. Thus we have a norm sequence

$$U_{h_2}U^{h_2}U^{t_2}.$$

###### Lemma 10.2.

There is a weak equivalence

$$U^{t_2}\mathrm{\Sigma }(\underset{𝑗}{\text{holim }}U_{h_2}S^{j\alpha })$$

###### Proof.

For each $`j`$ we have a cofibration sequence

$$(US^{(j1)\alpha })^{t_2}(US^{j\alpha })^{t_2}(U(_2)_+S^j)^{t_2}.$$

The spectrum $`U(_2)_+S^j`$ is the result of inducing $`US^j`$ up to a free $`_2`$-spectrum. Consequently, its Tate spectrum is weakly contractible (\[W-W, 2.6\]). It follows that the map $`(US^{(j1)\alpha })^{t_2}(US^{j\alpha })^{t_2}`$ is a weak equivalence for all $`j`$.

Next, consider for each $`j`$ the norm sequence

$$U_{h_2}S^{j\alpha }(US^{j\alpha })^{h_2}(US^{j\alpha })^{t_2}.$$

If we take the homotopy inverse limit over $`j`$, the middle term becomes weakly contractible (since homotopy fixed points commute with homotopy inverse limits and the maps $`US^{j\alpha }US^{(j+1)\alpha }`$ are unequivariantly null homotopic). The maps indexing the homotopy limit of the third term are all weak equivalences, so the homotopy limit in this case is just $`U^{t_2}=(US^0)^{t_2}`$. The result is now immediate. ∎

Let

$$u_f^{\mathrm{}}(E):WE\mathrm{\Sigma }^1(VE^2)^{t_2}$$

be the natural transformation resulting from Lemmas 10.1 and 10.2.

###### Corollary 10.3.

Up to homotopy, we have a factorization

where $``$ is the connecting map in the norm sequence.

###### Lemma 10.4.

When restricted to finite spectra, the functor

$$E\mathrm{\Sigma }^1(VE^2)^{t_2}$$

has degree $`1`$.

###### Proof.

(Sketch). There is a cofibration sequence

$$\mathrm{\Sigma }^1(VE^2)^{t_2}V_{h_2}E^2(VE^2)^{h_2}.$$

All functors in this sequence have degree $`2`$. The functor in the middle is homogeneous of degree two, and hence, its linearization vanishes. A direct calculation which we omit shows that the homogeneous degree two part of $`E(VE^2)^{h_2}`$ coincides with $`EV_{h_2}E^2`$ (this is where the assumption that $`E`$ ranges over finite spectra is used, for I don’t know whether the functor $`E(VE^2)^{h_2}`$ preserves filtered homotopy colimits). We infer that $`E\mathrm{\Sigma }^1(VE^2)^{t_2}`$ has trivial homogeneous degree two part; it therefore has degree $`1`$. ∎

A consequence of this lemma is that the natural transformation $`u_f^{\mathrm{}}`$ is determined by its value at $`S^0`$: i.e., for finite spectra $`E`$ we have $`u_f^{\mathrm{}}(E)u_f^{\mathrm{}}(S^0)\text{id}_E`$. Since $`u_f(E)u_f^{\mathrm{}}(E)`$, where $``$ is independent of $`f`$, it follows that the restriction of $`u_f`$ to finite spectra is determined by $`u_f^{\mathrm{}}(S^0)`$. As $`f`$ is determined by its restriction to finite spectra, we conclude

###### Corollary 10.5.

The functor $`f`$ is determined up to weak equivalence by the homotopy class

$$[u_f^{\mathrm{}}(S^0)][W,\mathrm{\Sigma }^1V^{t_2}].$$

###### Remark 10.6.

The above analysis yields an alternative proof of Theorem D. The point is that away from 2, the connecting map $``$ in the norm sequence is null homotopic.

We now apply 10.5 in the case of the functor $`ES_2(E)`$. Here, $`W=S^0`$ and $`V=S^1`$ (with trivial action). Furthermore, the natural transformation $`u_f(E)`$ is just $`\delta _E`$.

We see that $`\delta _E`$ is determined by a certain homotopy class

$$\delta ^{\mathrm{}}[S^0,(S^0)^{t_2}].$$

With slightly more care, it is not difficult to give an explicit construction: let map $`\iota :S^0(S^0)^{h_2}`$ be adjoint to the evident unit map $`(B_2)_+S^0`$. Then $`\delta ^{\mathrm{}}`$ represented by the composite

$$S^0\stackrel{𝜄}{}(S^0)^{h_2}(S^0)^{t_2}.$$

According to the Segal conjecture \[C\], $`\pi _0((S^0)^{t_2})`$ is isomorphic to the $`2`$-adic integers. The class $`\delta ^{\mathrm{}}`$ represents a topological generator of this group.


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# References

530.145+004.387

Classical and Quantum Problems for Quanputers  111The results of this paper where presented on the Workshop
”Quantum physics and communication ”, Dubna, May 16-17, 2002.

Nugzar Makhaldiani

Laboratory of Information Technologies

Joint Institute for Nuclear Research

Dubna, Moscow Region, Russia

e-mail address: mnv@jinr.ru

## Abstract

Quantum computers are considered as a part of the family of the reversible, lineary-extended, dynamical systems (Quanputers). For classical problems an operational reformulation is given. A universal algorithm for the solving of classical and quantum problems on quanputers is formulated.

1. Time-reversible classical discrete dynamical systems (computers) and corresponding (quantum) extraparts

The contemporary digital computer and its logical elements can be considered as a spatial type of discrete dynamical systems:

$`s_n(k+1)=\varphi _n(s(k)),`$ (1)
where $`s_n(k),1nN(k),`$ is the state vector of the system at the discrete time step $`k`$. Note that, with time steps $`k,`$ can change not only value, but also the dimension $`N(k),`$ of the state vector $`s(k)`$.

Definition. We assume that the system (1) is time-reversible if we can define the reverse dynamical system

$`s_n(k)=\varphi _n^1(s(k+1)).`$ (2)
In this case the following matrix

$`M_{nm}={\displaystyle \frac{\varphi _n(s(k))}{s_m(k)}},`$ (3)
is regular, i.e. has an inverse. If the matrix is not regular, this is the case, for example, when $`N(k+1)N(k),`$ we have an irreversible dynamical system (usual digital computers and/or corresponding irreversible gates).

Let us consider an extension of the dynamical system (1) given by the following action functional:

$`A={\displaystyle \underset{kn}{}}l_n(k)(s_n(k+1)\varphi _n(s(k)))`$ (4)
and corresponding motion equations

$`s_n(k+1)=\varphi _n(s(k))={\displaystyle \frac{H}{l_n(k)}},`$ (5)
$`l_n(k1)=l_m(k){\displaystyle \frac{\varphi _m(s(k))}{s_n(k)}}=l_m(k)M_{mn}(s(k))={\displaystyle \frac{H}{s_n(k)}},`$ (6)
where $`H=_{kn}l_n(k)\varphi _n(s(k))`$, is discrete Hamiltonian. In the regular case, we put this system in an explicit form

$`s_n(k+1)=\varphi _n(s(k)),`$ (7)
$`l_n(k+1)=l_m(k)M_{mn}^1(s(k+1)).`$ (8)
From this system it is obvious that, when the initial value $`l_n(k_0)`$ is given, the evolution of vector $`l(k)`$ is defined by evolution of the state vector $`s(k).`$ So we have the following

Theorem. The regular dynamical system (1) and the extended system (7) are equivalent.

Note that the corresponding statement from  concerning to the continual time dynamical systems is not so transparent as our statement for discrete time dynamical systems. In the continual time approximation, the discrete system (7) reduces to a corresponding continual one, :

$`\dot{x}_n(t)=v_n(x),\dot{p}_n(t)={\displaystyle \frac{v_m}{x_n}}p_m.`$ (9)
Indeed, let us change the dependent variables,

$`s_n(k)=x_n(t_k),l_n(k)=p_n(t_k),t_k=k\mathrm{\Delta }t,\mathrm{\Delta }t<<1.`$ (10)
Then the action functional (4) can be transformed as follows:

$`A`$ $`={\displaystyle \underset{kn}{}}p_n(t_k)(x_n(t_k+\mathrm{\Delta }t)\varphi _n(x(t_k)))`$ (13)
$`{\displaystyle \underset{kn}{}}\mathrm{\Delta }tp_n(t_k)(\dot{x}_n(t_k)v_n(x(t_k)))`$
$`{\displaystyle 𝑑tp_n(t)(\dot{x}_nv_n(x))},v_n(t_k)=(\varphi _n(x(t_k))x_n(t_k))/\mathrm{\Delta }t,`$
and corresponding motion equations are presented by the system (9). The equation of motion for $`l_n(k)`$ is linear and has an important property that a linear superpositions of the solutions are also solutions.

Statement. Any time-reversible dynamical system (e.g. a time-reversible computer) can be extended by a corresponding linear dynamical system (quantum-like processor) which is controlled by the dynamical system and which, due to the superposition and entanglement properties, has a huge computational power.

Note that  for the following (infinite - dimensional) equation of motion

$`iV_t=\mathrm{\Delta }V{\displaystyle \frac{1}{2}}V^2,`$ (14)
the corresponding linear subsystem is nothing but the familiar Schr$`\ddot{o}`$dinger equation,

$`i\psi _t=\mathrm{\Delta }\psi +V\psi .`$ (15)
The extended system

$`iV_t=\mathrm{}V{\displaystyle \frac{1}{2}}V^2,`$ (16)
$`i\psi _t=\mathrm{\Delta }\psi +V\psi ,`$ (17)
is mathematically and physically complete, as potential is not an external function, but is part of the Hamiltonian dynamical model.

2. Solution of the Quantum problems with the Methods of the Quantum Computing

The standard classical computing - the technology of storing and transforming information, is based on the classical physical theory and can be universally divided (factored) on the hardware (body) - memory and processor of computers, and software (soul) - algorithms and programs of users, parts. Under the Quantum computing (Quanputing, QC), we mean an extension of the classical computing by hardware which is described by the (corresponding extention of the) quantum theory - quantum registers and quantum logical units and corresponding software. In Quanputing the factorization on hardware and software is only an approximate concept, $`<SH>=<S><H>+O(\mathrm{}).`$

A quantum system can be described by corresponding Schr$`\ddot{o}`$dinger equation

$`i\mathrm{}{\displaystyle \frac{d}{dt}}|\psi >=\widehat{H}|\psi >,`$ (18)
where $`|\psi >=|\psi (t)>`$ is the state vector from the state Hilbert space and $`\widehat{H}=H(\widehat{p},\widehat{x})`$ is an operator-Hamiltonian. In the case of one nonrelativistic particle the operator is $`\widehat{H}=\widehat{T}+\widehat{V}=\frac{\widehat{p}^2}{2m}+V(\widehat{x}).`$ The fundamental bracket is $`[\widehat{x},\widehat{p}]=\widehat{x}\widehat{p}\widehat{p}\widehat{x}=i\mathrm{}.`$ The configuration space form of Eq. (18) is

$`i\mathrm{}{\displaystyle \frac{\psi (x,t)}{t}}=\widehat{H}\psi (x,t),`$ (19)
where $`\psi (x,t)=<x|\psi (t)>`$ and Hamiltonian is $`\widehat{H}=\frac{\mathrm{}^2}{2m}\frac{d^2}{dx^2}+V(x).`$

The formal solution of Eq. (18) is

$`|\psi (t)>=U(t)|\psi _0>,U(t)=exp({\displaystyle \frac{i}{\mathrm{}}}t\widehat{H}).`$ (20)
The main steps maid in QC are the following:

$`U(t)=(U^{1/N})^N=(U_TU_V)^N+O(1/N),`$ (21)
where

$`U^{1/N}`$ $`=`$ $`exp(\theta \widehat{H})=exp(\theta \widehat{T})exp(\theta \widehat{V})+O(1/N^2)`$ (22)
$`=`$ $`U_TU_V+O(1/N^2),\theta ={\displaystyle \frac{i}{\mathrm{}}}\tau ,\tau ={\displaystyle \frac{t}{N}}.`$ (23)
Then, for a corresponding matrix element we have (see Appendix)

$`<x_{n+1}|U^{1/N}|x_n>exp(\theta ({\displaystyle \frac{m}{2}}({\displaystyle \frac{x_{n+1}x_n}{\tau }})^2V(x_n)))+O(1/N^2)`$ (24)
and

$`<x_{out}|U(t)|x_{in}>`$ $``$ $`{\displaystyle 𝑑x_1𝑑x_2\mathrm{}𝑑x_Nexp(\frac{i}{\mathrm{}}\tau \underset{n=0}{\overset{N}{}}(\frac{m}{2}(\frac{x_{n+1}x_n}{\tau })^2V(x_n)))}`$ (25)
$`+`$ $`O(1/N),|x_0>=|x_{in}>,<x_{N+1}|=<x_{out}|.`$ (26)
This finite dimensional integral representation of the matrix element is in the ground of the functional (continual) integral formulation of the quantum theory . Then we discretize the wave function in (19), (20) and impose the periodic boundary conditions  $`a_n(t)=\psi (x_n,t),a_{n+N}=a_n,x_n=n\mathrm{\Delta }x.`$

We store these amplitudes in a k-bit quantum register,

$`|\psi (t)>={\displaystyle \underset{n=1}{\overset{N}{}}}a_n|n>,N=2^k,`$ (27)
where $`|n>`$ is the basis state corresponding to the binary representation of the number n. The second factor in (22) in a coordinate representation corresponds to a diagonal unitary transformation of the quantum computer state $`\psi (x,t).`$ After Fourier transformation $`\psi (x,t)`$ into momentum-space representation $`\psi (p,t)`$ the first factor in (22) can be applied in the same way. Diagonal unitary transformations of the type $`|n>e^{iF(n)}|n>,`$ where $`F(n)`$ is a function of n, can be done  with the following succession of steps

$`|n>`$ $``$ $`|n,0>|n,F(n)>e^{iF(n)}|n,F(n)>e^{iF(n)}|n,0>`$ (28)
$``$ $`e^{iF(n)}|n>.`$ (29)
For n particles in d dimensions we need $`nd`$ quantum registers. If the potential $`V(\widehat{x}_1,\widehat{x}_2,\mathrm{})`$ couples different degrees of freedom, we need the diagonal unitary transformations acting on several registers,

$`|n_1,n_2,\mathrm{}>e^{iF(n_1,n_2,\mathrm{})}|n_1,n_2,\mathrm{}>.`$ (30)
The disctetized bosonic quantum fields are naturally defined on the space-time lattice, scalar (particles) fields - on lattice sites, vector particles - on lattice links, tensor particles - on lattice plackets, and so on (see e.g. ).

3. Solution of the Classical Problems with the Methods of the Quantum Computing

If we can solve the quantum problems with the quantum computers , we can probably solve the classical problems, too. With the paper  we started the investigation of computational hard problems of classical physics, such as the turbulent phase of hydrodynamics, with the methods of Quantum Computing (QC).

We usually formulate the classical problems as a Hamiltonian system of motion equations

$`\dot{q}_n={\displaystyle \frac{H}{p_n}},\dot{p}_n={\displaystyle \frac{H}{q_n}},1nN,`$ (31)
where $`q_n`$ are coordinates for the configuration space of the system and $`p_n`$ are corresponding momentum, $`H=H(q,p)`$ is a Hamiltonian function(al). The system (31) belongs to a more general class of the dynamical systems defined by the following system of motion equations

$`\dot{x}_n=v_n(x),1nN,`$ (32)
when N is an even integer and

$$v_n(x)=\epsilon _{nm}\frac{H(x)}{x_m}=\{x_n,H\},$$

with the fundamental canonical bracket $`\{x_n,x_m\}=\epsilon _{nm}.`$

Note that any dynamical system (32) can be extended to the following Hamiltonian form, (see, e.g.  and section 1 of this paper.)

$`\dot{x}_n=v_n(x)=\{x_n,H\},`$ (33)
$`\dot{p}_n={\displaystyle \frac{v_m}{x_n}}p_m=\{p_n,H\},`$ (34)
with Hamiltonian $`H(x,p)=_nv_n(x)p_n`$ and the fundamental canonical bracket $`\{x_n,p_m\}=\delta _{nm}.`$

For any observable $`\psi (x)`$ of the Hamiltonian dynamical system, we have the following motion equation:

$`\dot{\psi }(x)=\{\psi ,H_1\}=\epsilon _{nm}{\displaystyle \frac{\psi }{x_n}}{\displaystyle \frac{H_1}{x_m}}=\widehat{L}\psi ={\displaystyle \frac{i}{\mathrm{}}}\widehat{H}_2\psi ,`$ (35)
where

$`\widehat{L}=v_n(x){\displaystyle \frac{}{x_n}}={\displaystyle \frac{i}{\mathrm{}}}v_n(x)\widehat{p}_n={\displaystyle \frac{i}{\mathrm{}}}\widehat{H}_2,`$ (36)
$`H_1`$ is the classical Hamiltonian function and $`\widehat{H}_2`$ is the quantum Hamiltonian operator. Now we can consider Eq.(35) as a classical Hamiltonian one with Hamiltonian $`H_1`$ or as a quantum Schr$`\ddot{o}`$dinger equation with the Hamiltonian operator $`\widehat{H}_2.`$ For the second form we can apply the QC methods developed in the previous section.

Note that the operator $`\widehat{H}_2`$ is a self-adjoint operator, $`\widehat{H}_2^+=\widehat{H}_2,`$ when $`divv=0`$ and we can put it in the second-order form with respect to the momentum operator,

$`\widehat{H}_2=v_k\widehat{p}_k={\displaystyle \frac{1}{2}}(\widehat{p}_kv_k+v_k\widehat{p}_k)={\displaystyle \frac{i}{\mathrm{}}}[\widehat{p}_k^2,u_k],`$ (37)
where

$`u_k={\displaystyle \frac{1}{2}}{\displaystyle \stackrel{x_k}{}}𝑑x_kv_k,1kd,d2.`$ (38)
The formal solution (20) is

$$|\psi (t)>=U(t)|\psi _0>=exp(\frac{i}{\mathrm{}}t\widehat{H}_2)|\psi _0>=exp(\frac{t}{\mathrm{}^2}[\widehat{u},\widehat{p}^2])|\psi _0>$$

and for the factorized form (21) we have

$`U^{1/N}=exp(\tau ^2[\widehat{u},\widehat{p}^2])={\displaystyle \underset{k=1}{\overset{d}{}}}e^{i\tau \widehat{u}_k}e^{i\tau \widehat{p}_k^2}e^{i\tau \widehat{u}_k}e^{i\tau \widehat{p}_k^2}+O(\tau ^3),`$ (39)
where

$$\tau =\pm (\frac{t}{\mathrm{}^2N})^{1/2}.$$

Now we are ready to apply the formalism of the QC considered in Sec.3 to the classical problems. In the extended version of this paper, we describe some computational problems from different parts of physics. The Hamiltonian dynamics of heavy-particle accelerators; the standard model of condensed state physics (Hubbard model); the 123 standard model of high-energy (elementary particles) physics; unified string and M-theory models.

4. Conclusions and perspectives

In this paper we constructed an algorithm for solving difficult computational(hard) problems of classical and quantum physics on quanputers. Note that before we get a real quanputer, we can simulate it on classical computers . It is interesting to investigate the functional integral formulation of the classical theory based on a discrete representation (39). Then, Dirac’s equation for electron

$`i\mathrm{}_t\psi (t,x)=(\alpha _n\widehat{p}_n+\beta m)\psi (t,x)=({\displaystyle \frac{i}{\mathrm{}}}[\widehat{p}_n^2,u_n]+\beta m)\psi ,`$ (40)
(where $`u_n=\frac{1}{2}𝑑x_n\alpha _n=\alpha _nx_n/2,`$) for $`m=0,`$ has the form similar to the classical equation (35), (with $`\alpha _n`$ as velocity matrices) and the Planck’s constant $`\mathrm{}`$ can be cancelled. This is an inverse to the consideration made when we try to find quantum (photon) nature of the (classical) Maxwell’s equations and put them in the ”Dirac’s form”, (see e.g. ). It is possible to formulate similar equations for high spin massless fields. It is interesting to find a corresponding classical interpretation of these equations and make observable predictions. The effects of two of them (the photon’s and graviton’s) are obvious in everyday life. What about the effects of other massless particles?

Appendix

1. In the main text of this paper we used the following relation:

$`e^{\epsilon A}e^{\epsilon B}e^{\epsilon A}e^{\epsilon B}=e^{\epsilon ^2[A,B]}+O(\epsilon ^3).`$ (41)
The relation

$`(1+\epsilon A)(1+\epsilon B)(1+\epsilon A)^1(1+\epsilon B)^1=(1+\epsilon ^2[A,B])+O(\epsilon ^3),`$ (42)
may also be useful.

2. Coordinate and momentum state vectors are correspondingly $`|x>`$ and $`|p>,\widehat{x}|x>=x|x>,\widehat{p}|p>=p|p>,`$

$`<p|x>=\psi _x(p)={\displaystyle \frac{1}{\sqrt{2\pi \mathrm{}}}}exp({\displaystyle \frac{i}{\mathrm{}}}px),\widehat{x}\psi _x(p)=i\mathrm{}{\displaystyle \frac{}{p}}\psi _x(p)=x\psi _x(p),`$ (43)
$`<x|y>={\displaystyle 𝑑p<x|p><p|y>}={\displaystyle \frac{dp}{2\pi \mathrm{}}exp(\frac{i}{\mathrm{}}p(xy))}`$ (44)
$`=\delta (xy).`$ (45)
Now we calculate the following matrix element

$`<x_{n+1}|exp(a\widehat{p}^2)|x_n>`$ $`=`$ $`{\displaystyle d^Dp<x_{n+1}|p><p|x_n>exp(ap^2)}`$ (46)
$`={\displaystyle \frac{d^Dp}{(2\pi \mathrm{})^D}exp(i\frac{p(x_{n+1}x_n)}{\mathrm{}}ap^2)}`$ $`=`$ $`{\displaystyle \frac{A^D}{(2\pi \mathrm{})^D}}exp({\displaystyle \frac{(x_{n+1}x_n)^2}{4a\mathrm{}^2}}),`$ (47)
where in the case of the quantum mechanics of the particle, (22)

$`a=i{\displaystyle \frac{t}{2m\mathrm{}N}}`$ (48)
and in the case of classical mechanics, (39)

$`a=i({\displaystyle \frac{t}{\mathrm{}^2N}})^{1/2},`$ (49)
$`A={\displaystyle 𝑑pexp(ap^2)}=\sqrt{{\displaystyle \frac{\pi }{a}}}.`$ (50)


text:
# The physics of stone skipping

## I Introduction

Nearly everyone has tried to throw a stone on a lake and count the number of bounces the stone was able to make. Of course the more, the better.skip  Our intuition gives us some empirical rules for the best throw: the best stones are flat and rather circular; one has to throw them rather fast and with a small angle with the water surface; a small kick is given with a finger to give the stone a spin. Of course these rules can be understood using the laws of physics: the crucial part of the motion is the collisional process of the stone with the water surface. The water surface exerts a reaction (lift) force on the stone, allowing it to rebound. This process is quite complex because it involves the description of the flow around the immersed stone.Stong ; Crane  Some energy is also dissipated during a collision, so that after a few rebounds, the initial kinetic energy of the stone is fully dissipated and the stone sinks.

The purpose of this paper is to propose a simplified description of the bouncing process of a stone on water, in order to estimate the maximum number of bounces performed by the stone. This problem provides an entertaining exercise for undergraduate students, with simple explanations for empirical laws that almost everyone has experienced.

## II Basic Assumptions

Consider a flat stone, with a small thickness and a mass $`M`$. The stone is thrown over a flat water surface. The angle between the stone surface and the water plane is $`\theta `$. A schematic view of the collisional process is shown in Fig. 1. The velocity $`V`$ is assumed to lie in a symmetry plane of the stone (the plane of the paper). The difficult part of the problem is, of course, to model the reaction force due to the water, which results from the flow around the stone during the stone-water contact. It is not the aim of this paper to give a detailed description of the fluid flow around the colliding stone. Rather I shall use a simplified description of the force keeping only the main ingredients of the problem. First, the velocity $`V`$ of the stone is expected to be (at least initially) the order of a few meters per second. For a stone with a characteristic size $`a`$ of the order of a few centimeters, the Reynolds number, defined as $`e=Va/\nu `$, with $`\nu `$ the kinematic viscosity ($`\nu 10^6\mathrm{m}^2\mathrm{s}^1`$ for water), is of order $`e10^5`$, that is, much larger than unity.Tritton  In this (inertial) regime, the force due to the water on the stone is expected on dimensional grounds to be quadratic in the velocity and proportional to the apparent surface of the moving object and the mass density of the fluid.Landau  Because the stone is only partially immersed in water during the collisional process, we expect the force to be proportional to the immersed surface (see Fig. 1). The force can be adequately decomposed into a component along the direction of the stone (that is, along $`𝐭`$, see Fig. 1) and a component perpendicular to it (that is, along $`𝐧`$). The latter corresponds to the lift component of the force, and the former corresponds to a friction component (of water along the object). I write the reaction force due to water, $`𝐅`$, as:

$$𝐅=\frac{1}{2}C_l\rho _wV^2S_{\mathrm{im}}𝐧+\frac{1}{2}C_f\rho _wV^2S_{\mathrm{im}}𝐭,$$
(1)
where $`C_l`$ and $`C_f`$ are the lift and friction coefficients, $`\rho _w`$ is the mass density of water, $`S_{\mathrm{im}}`$ is the area of the immersed surface, and $`𝐧`$ is the unit vector normal to the stone (see Fig. 1).

Note that in general, both $`C_l`$ and $`C_f`$ are functions of the tilt angle $`\theta `$ and incidence angle $`\beta `$, defined as the angle between velocity $`V`$ and the horizontal. In the simplified analysis I will assume that both $`C_l`$ and $`C_f`$ are constant and independent of tilt and incidence angles.Note  This assumption is not a strong one because ricochets are generally performed with a small tilt angle, $`\theta `$, and a small incidence angle, $`\beta `$. If one denotes the initial components of the incident velocity by $`V_{x0}`$ and $`V_{z0}`$ (parallel and perpendicular to the water surface, respectively), the latter assumption amounts to $`V_{z0}V_{x0}`$.

We expect the lift force to be maximum when the object is only partially immersed due to the lack of symmetry between the two sides of the stone. Therefore, if the object reaches a depth such that it becomes completely immersed, the lift force would be greatly diminished and would probably not be able to sustain the weight of the stone anymore. For simplicity, I will assume that the lift force vanishes for completely immersed objects. The model for the force in Eq. (1) is crude, but it is expected to capture the main physical ingredients of the stone-water interaction. It might fail for lower stone velocities or larger incidence angles, where a bulge of water could be created and affect the lift and friction forces on the stone.Stong  However, in this case it is expected that the stone will be strongly destabilized during the collision process and perform only a very small number of bounces. We will restrict ourselves to large initial velocities and small incidence angles, such that the number of bounces is sufficiently large.

## III Equations of motion

Consider the collisional process, that is, the time during which the stone is partially immersed in water. I will assume in this section that the incidence angle $`\theta `$ between the stone and the water surface is constant during the collisional process. The validity of this assumption is considered in detail in Sec. V. The origin of time, $`t=0`$, corresponds to the instant when the edge of the stone reaches the water surface. During the collisional process, the equations of motion for the center of mass velocity are

$`M{\displaystyle \frac{dV_x}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\rho _wV^2S_{\mathrm{im}}(C_l\mathrm{sin}\theta +C_f\mathrm{cos}\theta )`$ (2a)
$`M{\displaystyle \frac{dV_z}{dt}}`$ $`=`$ $`Mg+{\displaystyle \frac{1}{2}}\rho _wV^2S_{\mathrm{im}}(C_l\mathrm{cos}\theta C_f\mathrm{sin}\theta ),`$ (2b)
with $`V^2=V_x^2+V_z^2`$ and $`g`$ is the acceleration due to gravity. Note that in Eq. (2) the area $`S_{\mathrm{im}}`$ depends on the immersed depth and thus varies during the collisional process.

Equation (2) is non-linear due to the $`V^2`$ terms on the right-hand side, but also due to the dependence of the immersed area, $`S_{\mathrm{im}}`$, on the height $`z`$. However, we can propose a simple approximation scheme: the magnitude of the velocity, $`V`$, is not expected to be strongly affected by the collision process (as I shall show in Sec. VI). I thus make the approximation that $`V^2V_{x0}^2+V_{z0}^2V_{x0}^2`$ on the right-hand side of Eq. (2). The validity of this assumption requires a sufficiently high initial velocity, $`V_{x0}`$, and it might fail in the last few rebounds of a stone skip sequence.

With this approximation, Eq. (2b) decouples from Eq. (2a). I thus first focus the discussion on the equation for the height $`z`$, which is the height of the immersed edge (see Fig. 1). Note that the equation for $`z`$ is equivalent to the equation of the center of mass position, Eq. (2b) because $`\theta `$ is assumed to be constant (see Sec. V for a detailed discussion of this point). Hence, we may identify $`V_z`$ with $`dz/dt`$ and Eq. (2b) yields a closed equation for the height $`z`$.

## IV Collisional process

To solve Eq. (2b) we need to prescribe the $`z`$ dependence of the immersed area $`S_{\mathrm{im}}`$. This quantity depends on the precise shape of the stone. A natural choice is circular, which I will treat in Sec. IV.2. However, it is enlighting to first consider a square shape; this shape greatly simplifies the mathematics and already contains the basic mechanisms involved.

### IV.1 A Square Stone

In this case, the immersed area is simply $`S_{\mathrm{im}}=a|z|/\mathrm{sin}\theta `$ (see Fig. 1), with $`a`$ the length of one edge of the stone. The equation for $`z`$ thus becomes

$$M\frac{d^2z}{dt^2}=Mg\frac{1}{2}\rho _wV_{x0}^2C\frac{az}{\mathrm{sin}\theta },$$
(3)
where $`C=C_l\mathrm{cos}\theta C_f\mathrm{sin}\theta C_l`$, and I have used $`|z|=z`$ ($`z<0`$). We define the characteristic frequency $`\omega _0`$ as

$$\omega _0^2=\frac{C\rho _wV_{x0}^2a}{2M\mathrm{sin}\theta },$$
(4)
and rewrite Eq. (3) as

$$\frac{d^2z}{dt^2}+\omega _0^2z=g.$$
(5)
With the initial conditions at $`t=0`$ (first contact with water), $`z=0`$ and $`\dot{z}=V_{z0}<0`$, the solution of Eq. (5) is

$$z(t)=\frac{g}{\omega _0^2}+\frac{g}{\omega _0^2}\mathrm{cos}\omega _0t+\frac{V_{z0}}{\omega _0}\mathrm{sin}\omega _0t.$$
(6)
Equation (6) characterizes the collisional process of the stone with water. After a collision time $`t_{\mathrm{coll}}`$ defined by the condition $`z(t_{\mathrm{coll}})=0`$ ($`t_{\mathrm{coll}}2\pi /\omega _0`$), the stone emerges totally from the water surface. It is easy to show that the maximal depth attained by the stone during the collision is

$$|z_{\mathrm{max}}|=\frac{g}{\omega _0^2}\left[1+\sqrt{1+\left(\frac{\omega _0V_{z0}}{g}\right)^2}\right].$$
(7)
As discussed in Sec. I, the stone will rebound if it stays only partially immersed during the collision. The rebound condition can be written as $`|z_{\mathrm{max}}|<a\mathrm{sin}\theta `$. If we use Eqs. (7) and (4), this condition can be written after some straightforward calculations as

$$V_{x0}>V_c=\frac{\sqrt{\frac{4Mg}{C\rho _wa^2}}}{\sqrt{1\frac{2\mathrm{tan}^2\beta M}{a^3C\rho _w\mathrm{sin}\theta }}},$$
(8)
where the incidence angle $`\beta `$ is defined as $`V_{z0}/V_{x0}=\mathrm{tan}\beta `$. Therefore, we obtain a minimum critical velocity for skimming. Using the typical values, $`M=0.1`$ kg, $`L=0.1`$ m, $`C_lC_f1`$, $`\rho _w=1000`$ kg m-3, and $`\beta \theta 10^{}`$, we obtain $`V_c0.71\mathrm{m}\mathrm{s}^11`$ m s-1.

The physical meaning of this condition is clear: it simply expresses the fact that the lift force $`\frac{1}{2}C\rho _wV^2a^2`$ has to balance the weight of the stone $`Mg`$ in order for it to bounce.

### IV.2 A circular stone

For a circular stone, the immersed area is a more complex function of the height $`z`$, and is given in terms of the area of a truncated circle. A simple integral calculation yields

$$S_{\mathrm{im}}(s)=R^2[\mathrm{arccos}(1s/R)(1s/R)\sqrt{1(1s/R)^2}],$$
(9)
with $`s=|z|/\mathrm{sin}\theta `$ (the maximum immersed length) and $`R=a/2`$ as the radius of the stone.

The equation of motion for $`z`$, Eq. (2b), thus becomes non-linear. However, it is possible to describe (at least qualitatively) the collisional process and obtain the condition for the stone to bounce.

I first introduce dimensionless variables to simplify the calculations. The dimensionless height, $`\stackrel{~}{z}`$, time, $`\tau `$, and immersed area, $`𝒜`$, are defined as $`\stackrel{~}{z}=z/R\mathrm{sin}\theta `$, $`\tau =\omega _0t`$, and $`𝒜(\stackrel{~}{z})=S_{\mathrm{im}}/R^2`$. (The minus sign in $`\stackrel{~}{z}`$ is introduced for convenience.) If we use these variables, Eq. (2b), and $`V_z=dz/dt`$, we obtain

$$\frac{d^2\stackrel{~}{z}}{d\tau ^2}=\alpha \frac{1}{2}𝒜(\stackrel{~}{z}),$$
(10)
with $`\alpha =g/(R\omega _0^2\mathrm{sin}\theta )`$. Equation (10) is the equation of a particle (with unit mass) in the potential $`𝒱(\stackrel{~}{z})=(\frac{1}{2}𝒜(\stackrel{~}{z})\alpha )𝑑\stackrel{~}{z}`$. We can use standard techniques for mechanical systems to solve Eq. (10). In particular, Eq. (10) can be integrated once to give the “constant energy” condition

$$\frac{1}{2}\left(\frac{d\stackrel{~}{z}}{d\tau }\right)^2+𝒱(\stackrel{~}{z})=E,$$
(11)
where $`E`$ is the energy of the system and is given in terms of the initial conditions

$$E=\frac{1}{2}\left(\frac{d\stackrel{~}{z}}{d\tau }\right)^2|_{\tau =0}+𝒱(\stackrel{~}{z}=0)=\frac{1}{2}(V_{z0}/(R\omega _0\mathrm{sin}\theta ))^2.$$
(12)
The potential $`𝒱(\stackrel{~}{z})`$ can be calculated analytically using the expression for the immersed area $`S_{\mathrm{im}}`$ given in Eq. (9). An integral calculation gives

$`𝒱(\stackrel{~}{z})=`$ $`{\displaystyle \frac{1}{2}}(\sqrt{1(1\stackrel{~}{z})^2}[{\displaystyle \frac{2}{3}}+{\displaystyle \frac{1}{3}}(1\stackrel{~}{z})^2]`$ (13)
$`(1\stackrel{~}{z})\mathrm{arccos}(1\stackrel{~}{z}))\alpha \stackrel{~}{z}.`$
This potential is plotted in Fig. 2 as a function of $`\stackrel{~}{z}`$. As a consequence of the constant energy condition, Eq. (11), $`\stackrel{~}{z}`$ exhibits a turning point at a maximum depth defined by $`𝒱(\stackrel{~}{z}_{\mathrm{max}})=E`$.

Here again, the condition for the stone to bounce is that this maximum depth be reached before the stone is fully immersed, that is, $`|z_{\mathrm{max}}|<2R\mathrm{sin}\theta `$. In terms of dimensionless variables, we obtain the condition: $`\stackrel{~}{z}_{\mathrm{max}}<2`$, with $`\stackrel{~}{z}_{\mathrm{max}}`$ defined by $`𝒱(\stackrel{~}{z}_{\mathrm{max}})=\frac{1}{2}(V_{z0}/(R\omega _0\mathrm{sin}\theta ))^2`$. This condition can be explicitly solved. Let me introduce $`\stackrel{~}{z}_0`$ such that $`d𝒱/d\stackrel{~}{z}=0`$ at $`\stackrel{~}{z}=\stackrel{~}{z}_0`$: $`𝒱(\stackrel{~}{z})`$ is a monotonically increasing function of $`\stackrel{~}{z}`$ for $`\stackrel{~}{z}>\stackrel{~}{z}_0`$. Now it is easy to show that $`\stackrel{~}{z}_{\mathrm{max}}>z_0`$ (because $`𝒱(\stackrel{~}{z}_{\mathrm{max}})>0`$ and $`𝒱(\stackrel{~}{z}_0)<0`$), and the condition $`\stackrel{~}{z}_{\mathrm{max}}<2`$ is therefore equivalent to $`𝒱(\stackrel{~}{z}_{\mathrm{max}})<𝒱(2)=\frac{\pi }{2}2\alpha `$, that is, $`\frac{1}{2}(V_{z0}/(R\omega _0\mathrm{sin}\theta ))^2<\frac{\pi }{2}2g/(R\omega _0^2\mathrm{sin}\theta )`$. Then the condition for skimming can be rewritten (recalling that $`V_{z0}/V_{x0}=\mathrm{tan}\beta `$)

$$V_{x0}>V_c=\frac{\sqrt{\frac{16Mg}{\pi C\rho _wa^2}}}{\sqrt{1\frac{8M\mathrm{tan}^2\beta }{\pi a^3C\rho _w\mathrm{sin}\theta }}}.$$
(14)
Up to (slightly different) numerical factors this condition is the same as in Eq. (8) for a square stone. Note moreover, that the reasoning used for the potential $`𝒱`$ is quite general and can be applied to the square shape as well. This reasoning yields the same condition as Eq. (8) in this case.

Note also that for the circular stone, a simplified analysis of the motion could have been performed. First if $`\stackrel{~}{z}`$ remains small during the bounce of the stone, a small $`\stackrel{~}{z}`$ expansion of $`𝒱(\stackrel{~}{z})`$ is possible, yielding $`𝒱(\stackrel{~}{z})=4\sqrt{2}/15\stackrel{~}{z}^{5/2}\alpha \stackrel{~}{z}`$ (corresponding to a parabolic approximation for the shape of the stone near its edge). Moreover, we remark that for small $`V_{z0}`$, the energy $`E`$ goes to zero, so that $`\stackrel{~}{z}_{\mathrm{max}}`$ is defined in this case by $`𝒱(\stackrel{~}{z}_{\mathrm{max}})=0`$. If also use the previous approximation, we obtain $`\stackrel{~}{z}_{\mathrm{max}}=(15\alpha /4\sqrt{2})^{2/3}`$. The condition for the stone to bounce, $`\stackrel{~}{z}_{\mathrm{max}}<2`$, therefore yields $`\alpha <16/15`$. In terms of $`V_{x0}`$, this condition gives again a minimum critical velocity for skimming, defined as $`V_c=\sqrt{\zeta Mg/C\rho _wa^2}`$ with $`\zeta =15/43.75`$. This result is thus close to the “exact” condition found in Eq. (14) for the $`V_{z0}=0`$ case.

### IV.3 Energy Dissipation

I have so far described the rebound of the stone by analyzing its vertical motion. This analysis gave a minimum velocity for skimming which results from the balance between the weight of the stone and the lift of the force due to water. However, some energy is dissipated during the collision due to the “friction” contribution of the force (the component along $`x`$). This mechanism of dissipation leads to another minimum velocity condition, in terms of the balance between dissipation and initial kinetic energy. Only a qualitative description of the dissipation is given here.

As shown by Eq. (2), the component $`F_x`$ of the reaction force in the $`x`$ direction (parallel to the water surface) will decrease the velocity of the stone. Then after a few bounces, the condition for the stone to bounce, Eq. (8) or Eq. (14), will no longer be satisfied and the stone will stop. It is possible to estimate the decrease in the $`x`$-component of the velocity using the equation for the center of mass position, Eq. (2a). If we multiply both sides of Eq. (2a) by $`V_x`$ and integrate over a collision time, we obtain the decrease in the kinetic energy in the $`x`$ direction in terms of the work of the reaction force:

$$𝒲\frac{1}{2}MV_{xf}^2\frac{1}{2}MV_{x0}^2=_0^{t_{\mathrm{coll}}}F_x(t)V_x(t)𝑑t,$$
(15)
where $`V_{x0}`$ and $`V_{xf}`$ are the $`x`$-components of the velocity before and after the collision, $`t_{\mathrm{coll}}`$ is the collision time, and $`F_x=\frac{1}{2}\stackrel{~}{C}\rho _wV_x^2S_{\mathrm{im}}`$ is the $`x`$-component of the reaction force, with $`\stackrel{~}{C}=C_l\mathrm{sin}\theta +C_f\mathrm{cos}\theta `$.

A rough estimate of the right-hand side of Eq. (15) is

$$_0^{t_{\mathrm{coll}}}F_x(t)V_x(t)𝑑tV_{x0}_0^{t_{\mathrm{coll}}}F_x(t)𝑑t.$$
(16)
Now we have the simple relation $`F_x(t)=\mu F_z(t)`$, with $`\mu =\stackrel{~}{C}/C`$ (see Eq. (1)). Moreover, it is expected that the average vertical force during a collision, $`F_z(t)=t_{\mathrm{coll}}^1_0^{t_{\mathrm{coll}}}F_z(t)𝑑t`$, is the order of the weight of the stone, $`Mg`$. This point can be explicitly verified for the square stone case, using the expression of the force $`F_z`$ in terms of the height $`z(t)`$ and Eq. (6). The final result is $`F_x(t)\mu Mg`$.Note2 . Moreover, as shown in the above (and in particular for the square stone, although the results remain qualitatively valid for the circular one), the collision time is given approximatively by $`t_{\mathrm{coll}}2\pi /\omega _0`$. We eventually find that the loss in kinetic energy in Eq. (15) is approximatively given by

$$𝒲\mu MgV_{x0}\frac{2\pi }{\omega _0}=\mu Mg\mathrm{},$$
(17)
where $`\mathrm{}`$ is defined as

$$\mathrm{}=V_{x0}\frac{2\pi }{\omega _0}=2\pi \sqrt{\frac{2M\mathrm{sin}\theta }{C\rho _wa}}.$$
(18)
The quantity $`\mathrm{}=V_xt_{\mathrm{coll}}`$ is the distance along $`x`$ traversed by the stone during a collision. If the energy loss $`𝒲`$ is larger than the initial kinetic energy, the stone would be stopped during the collision. Using Eq. (15), this condition can be written explicitly as $`\frac{1}{2}MV_{x0}^2>|𝒲|=\mu Mg\mathrm{}`$. We deduce that the initial velocity should be larger than the minimum velocity $`V_c`$ in order to perform at least one bounce, that is,

$$V_{x0}>V_c=\sqrt{2\mu g\mathrm{}}.$$
(19)
If we use the same numerical values as in the previous paragraph, we obtain $`\mu =1.4`$, $`\mathrm{}=13`$ cm, so that $`V_c2`$ m s-1. This criterion is more restrictive than the previous one, Eq. (14). I thus consider in the following that Eq. (19) is the criterion for the stone to skim over water.

## V Why give the stone a spin?

The previous calculations assumed a constant angle $`\theta `$. It is obvious that the rebound of the stone is optimized when $`\theta `$ is small and positive (see, for example, the value of the force constant $`C=C_l\mathrm{cos}\theta C_f\mathrm{sin}\theta `$ which decreases when $`\theta `$ increases). Now, if after a collision, the stone is put in rotation around the $`y`$-axis (see Fig. 1), that is, $`\dot{\theta }0`$, its orientation would change by an appreciable amount during free flight: the incidence angle $`\theta `$ for the next collision has little chance to still be in a favorable situation. The stone performs, say, at most one or two more collisions. There is therefore a need for a stabilizing angular motion. This is the role of the spin of the stone.

Let us denote $`\dot{\varphi }_0`$ as the rotational velocity of the stone around the symmetry axis parallel to $`𝐧`$ in Fig. 1. I neglect in the following any frictional torque on the stone (associated with rotational motion). During the collision, the reaction force due to the water is applied only to the immersed part of the stone and results in a torque applied on the stone. For simplicity, I consider only the lift part of the force. Its contribution to the torque (calculated at the center $`O`$ of the stone) can be readily calculated as $`_{\mathrm{lift}}=OPF_{\mathrm{lift}}𝐞_y`$, where $`𝐞_y`$ is the unit vector in the $`y`$ direction in Fig. 1 and $`P`$, the point of application of the lift force, is located at the center of mass of the immersed area. This torque is in the $`y`$ direction and will eventually affect the angular motion along $`\theta `$. However a spin motion around $`𝐧`$ induces a stabilizing torque: this is the well-known gyroscopic effect.Goldstein  The derivation of the equation of motion of the rotating object (the Euler equations) is a classic problem and is treated in standard mechanics textbooks (see for example, Ref. Goldstein, ). On the basis of these equations, it is possible to derive the stabilizing gyroscopic effect. This derivation is briefly summarized in the Appendix.

In our case, the equation for the angle $`\theta `$ can be written as

$$\ddot{\theta }+\omega ^2(\theta \theta _0)=\frac{_\theta }{J_1},$$
(20)
where $`\omega =[(J_0J_1)/J_1]\dot{\varphi }_0`$, $`\dot{\varphi }_0`$ is the initial spin angular velocity (in the $`𝐧`$ direction), and $`J_0`$ and $`J_1`$ are moments of inertia in the $`𝐧`$ and $`𝐭`$ directions, respectively; $`\theta _0`$ is the initial tilt angle and $`_\theta =OPF_{\mathrm{lift}}`$ is the projection of the torque due to the water flow in the $`y`$ direction. Equation (20) shows that in the absence of spin motion, $`\dot{\varphi }_0=0`$, the torque due to the lift force will initiate rotational motion of the stone in the $`\theta `$ direction. As discussed above, the corresponding situation is unstable. On the other hand, spin motion induces a stabilizing torque that can maintain $`\theta `$ around its initial value. The effect of the torque can be neglected if, after a collision with the water, the maximum amplitude of the motion of the angle $`\theta `$ is small: $`\delta \theta =[\theta \theta _0]_{\mathrm{max}}1`$. If we use Eq. (20), an estimate of $`\delta \theta `$ can be obtained by balancing the last two terms in Eq. (20), yielding $`\delta \theta _\theta /(J_1\omega ^2)`$ (note that up to numerical factors ($`J_0J_1)/J_1)1`$ and $`J_1MR^2`$, with $`R`$ the radius of the stone). The order of magnitude of $`_\theta `$ can be obtained using the results of Sec. IV.3. The average vertical force acting on the stone has been found to be the order of the weight of the stone (see the discussion after Eq. (16)): $`F_z(t)Mg`$. If we take $`OPR`$, we obtain the simple result $`_\theta MgR`$. The estimate for $`\delta \theta `$ follows directly as $`\delta \theta g/(R\omega ^2)`$. Therefore, the condition for $`\theta `$ to remain approximately constant, $`\delta \theta 1`$, is

$$\dot{\varphi }_0\omega \sqrt{\frac{g}{R}}.$$
(21)
For a stone with a diameter of 10 cm, Eq. (21) gives $`\dot{\varphi }_014`$ s-1, corresponding to a rotational frequency larger than a few revolutions per second ($`2`$ Hz). This condition is easily fullfilled in practice and corresponds approximately to what we would expect intuitively for a successful throw. Note that the condition (21) is independent of the center of mass velocity of the stone $`V`$.

## VI An estimate for the maximum number of bounces

The estimation of the maximum number of bounces is the most difficult and tentative part of the analysis because many factors can in principle slow down or destabilize the stone, some of which are extremely difficult to model (such as irregularities of the water surface and the wind). We shall assume the idealized situation described above (perfect surface, no wind, idealized reaction force) and focus on two specific factors, which appear, at least intuitively, as natural candidates for stopping the stone.

### VI.1 Slow down of the stone

As I have discussed in Sec. IV.3, energy is dissipated during a collision and the $`x`$ component of the velocity of the stone will decrease during each collision: after a few collisions, all the initial kinetic energy will be dissipated. This process can be easily formulated.

I consider a succession of $`N`$ collisions. Between two collisions, the motion is parabolic (wind and air friction are neglected) and the initial $`x`$ component of the velocity at the next collision is equal to the final $`x`$ component of the velocity at the end of the previous collision. The important point to note is that the energy loss during one collision, Eq. (17), is independent of the velocity $`V_{x0}`$ before the collision. Therefore, the velocity of the stone after $`N`$ collisions obeys the relation

$$\frac{1}{2}MV_x^2[N]\frac{1}{2}MV_x^2[0]=N\mu Mg\mathrm{},$$
(22)
so that the stone will be stopped at a collision number $`N_c`$ such that the total energy loss is larger than the initial kinetic energy (similar to the argument leading to the critical velocity for skimming, $`V_c`$, in Eq. (19)). This criterion corresponds to $`V_x^2[N_c]=0`$ in Eq. (22), and $`N_c`$ is given accordingly by

$$N_c=\frac{V_x^2[0]}{2g\mu \mathrm{}}.$$
(23)
If we use the same typical values as before ($`M=0.1`$ kg, $`a=0.1`$ m, $`C_lC_f1`$, $`\rho _w=1000`$ m-3, $`\beta \theta 10^{}`$, we obtain $`\mu 1.4`$ and $`\mathrm{}13`$ cm. We then find $`N_c6`$ for the initial velocity $`V_{x0}=5`$ m s-1, $`N_c17`$ for $`V_{x0}=8`$ m s-1, and $`N_c38`$ for $`V_{x0}=12`$ m s-1. The latter number of bounces corresponds to the world record.skip

It is interesting to calculate the distance between two successive collisions. As noted, the motion of the stone is parabolic out of the water: the position {$`X,Z`$} of the particle is given by $`X(t)=V_xt`$, $`Z(t)=\frac{1}{2}gt^2+|V_z|t`$. The next collision will occur at a distance $`\mathrm{\Delta }X=2V_x|V_z|/g`$. The dependence of $`V_x`$ on the number of collisions $`N`$ is given by Eq. (22). On the other hand, $`V_z`$ does not depend on the number of collisions because the stone rebounds “elastically” in the $`z`$ direction, as follows from the analysis of the collisional process in Sec. IV (see, for example, the conservation of the energy $`E`$ during the collision discussed for the circular stone). If we use Eq. (22), we obtain the simple result

$$\mathrm{\Delta }X[N]=\mathrm{\Delta }X_0\sqrt{1\frac{N}{N_c}},$$
(24)
where $`\mathrm{\Delta }X_0=2V_{x0}|V_{z0}|/g`$. Note that $`\mathrm{\Delta }X_0`$ is approximately equal to the distance between the two first ricochets, $`\mathrm{\Delta }X[N=1]`$, when $`N_c1`$. For $`V_{x0}=8`$ m s-1, we obtain $`\mathrm{\Delta }X_02.25`$ m.

Equation (24) for $`\mathrm{\Delta }X[N]`$ is plotted in Fig. 3. We remark that the decrease in the distance between two successive ricochets is first rather slow ($`\mathrm{\Delta }X[N]\mathrm{\Delta }X_0(1\frac{N}{2N_c})`$ for $`NN_c`$, see Eq. (24)), but strongly accelerates for the last collisions when $`NN_c`$, due to the square root variation of $`\mathrm{\Delta }X[N]`$ close to $`N_c`$. This result is in agreement with observation. Such an effect is known to specialists of stone-skipping as “pitty-pat.”skip

### VI.2 Angular destabilization

However, there is another possible destabilizing mechanism in the collision process. As was discussed in Sec. V, the rotational stability of the stone is crucial in the collisional process. A criterion for stability has been found in the form of a minimum spin velocity of the stone. However, each collision will perturb the rotational motion and the sum of all these effects can eventually bypass the stability condition. This argument can be easily formulated. As shown above, the amplitude of the angular motion of $`\theta `$ is $`\delta \theta g/(R\omega ^2)`$, with $`\omega \dot{\varphi }_0`$, the (constant) spin velocity of the stone. Now assume that the destabilizing effects add, a reasonable assumption. Then, after $`N`$ collisions we expect that $`\mathrm{\Delta }_N\theta N\delta \theta `$. The stone is completely destabilized for a collision number $`N_c`$ such that $`\mathrm{\Delta }_{N_c}\theta 1`$, yielding

$$N_c\frac{R\dot{\varphi }_0^2}{g}.$$
(25)
If we use the same numerical values as before, we obtain, for example, $`N_c5`$ for a initial spin velocity $`\varphi _0=5`$ rev/s and $`N_c=38`$ (the world recordskip ) for $`\varphi _0=14`$ rev/s. Note, however, that there is a quite large uncertainty of the numerical prefactors in the above estimate of $`N_c`$, and this estimate is merely qualitative and should not be taken literally.

## VII Discussion

At the level of our description, the maximum number of bounces results from the combination of the two previous mechanisms: slow down and angular destabilization. The maximum number of bounces is therefore given by the minimum of the two previous estimates, in Eqs. (23) and (25).

The estimate $`N_c^{\mathrm{sd}}`$ obtained in Eq. (23) from the slow down of the stone depends only (quadratically) on the initial velocity of the stone: in principle, a very large number of bounces could be reached by increasing the initial velocity of the stone. But on the other hand, the angular destabilization process results in a maximum value of $`N_c^{\mathrm{spin}}`$ which is independent of the initial velocity of the stone, as indicated by Eq. (25). This shows that even if the initial velocity of the stone is very large, that is, $`N_c^{\mathrm{sd}}1`$, the stone will be stopped by angular destabilization after $`N_c^{\mathrm{spin}}`$ bounces. In other words, the initial “kick” that puts the stone in rotational motion is a key factor for a good throw.

The results presented here are in agreement with our intuition for the conditions of a good throw. Some of the results are also in agreement with observations, for example, the acceleration of the number of collisions at the end of the throw (a phenomenon known as “pitty-pat” in stone skipping competitions skip ). Some easy checks of the assumptions underlying our calculations could be performed, even without any sophisticated apparatus. For example, taking pictures of the water surface after the ricochets would locate the positions of the collisions (because small waves are produced at the surface of water). A simple test of the variation of the distance between two collisions as a function of collision number, Eq. (24), would then be possible. A more ambitious project would be to design a “catapult,” allowing one to throw stones with a controlled translational and spin velocity (together with the incidence angle of the stone on water). A measurement of the maximum number of bounces performed for various throw parameters would allow us to check the assumptions underlying the present simple analysis and to determine some of the parameters involved in the description (such as $`\mu `$ and $`\mathrm{}`$). It would be also interesting to repeat the experiments reported in Ref. Stong,  using modern techniques (such as fast cameras), in order to image and analyze in particular the rebound process as a function of the throw parameters. Hopefully a better understanding of the mechanisms of stone skipping will allow someone to break the actual world record!

###### Acknowledgements.

I thank my son Léonard for his (numerous and always renewed) perplexing questions. I thank my colleagues from the physics laboratory of the ENS Lyon, and in particular Bernard Castaing and Thierry Dauxois, for their constant interest in discussing “simple” physics problems. I am grateful to the referee for pointing out Refs. Stong,  and Crane,  to me.

## Appendix A appendix

I briefly recall the derivation of Eq. (20), from the Euler equations described in Ref. Goldstein,  The latter are written as:Goldstein

$`I_1{\displaystyle \frac{d\omega _1}{dt}}\omega _2\omega _3(I_2I_3)`$ $`=N_1`$ (26a)
$`I_2{\displaystyle \frac{d\omega _2}{dt}}\omega _1\omega _3(I_3I_1)`$ $`=N_2`$ (26b)
$`I_3{\displaystyle \frac{d\omega _3}{dt}}\omega _1\omega _2(I_1I_2)`$ $`=N_3.`$ (26c)
In Eq. (26), $`I_\alpha `$, $`\omega _\alpha `$, and $`N_\alpha `$ ($`\alpha =1,2,3`$) are respectively the moment of inertia, angular velocity, and torque along the direction of a particular principal axis, denoted as $`\alpha `$. In our case, the direction 1 is taken along the axis perpendicular to the vectors $`𝐧`$ and $`𝐭`$ (the direction 1 is along the $`y`$ axis in Fig. 1), the direction 2 along $`𝐧`$ and the direction 3 along $`𝐭`$. We therefore have $`\omega _1=\dot{\theta }`$, and due to the symmetry of the circular stone, $`I_1=I_3J_1`$ and $`I_2J_0`$. Moreover, because only the lift component of the reaction force (along $`𝐧`$) is considered in the present analysis, we have $`N_1_\theta `$ and $`N_2=N_3=0`$.

Equation (26b) yields immediately that $`\dot{\omega }_2=0`$. We therefore have $`\omega _2=\dot{\varphi }_0`$, with $`\dot{\varphi }_0`$ the initial spin velocity. Equation (26c) can be therefore written as:

$$\frac{d\omega _3}{dt}=\frac{J_1J_0}{J_1}\dot{\varphi }_0\omega _1.$$
(27)
If we use $`\omega _1=\dot{\theta }`$, Eq. (27) can be integrated once to give

$$\omega _3=\frac{J_1J_0}{J_1}\dot{\varphi }_0(\theta \theta _0),$$
(28)
with $`\theta _0=\theta (t=0)`$, the initial tilt angle. The substitution of Eq. (28) into Eq. (26a) leads to Eq. (20).