format: text

text:
# 1 Introduction

## 1 Introduction

The minimal supersymmetric standard model (MSSM)  is one of the most attractive models beyond the standard model (SM). In the MSSM, two Higgs doublets are introduced in order to cancel the anomaly and to give the fermion masses. The introduction of the second Higgs doublet inevitably means that a charged Higgs boson is in the physical spectra. So, it is quite important to study effects of the charged Higgs boson.

Here, we study effects of the charged Higgs boson on the exclusive semi-tauonic $`B`$ decay, $`\overline{B}D\tau \overline{\nu }_\tau `$, in the MSSM. In our previous works , we calculated the decay rate of $`\overline{B}D\tau \overline{\nu }_\tau `$ including the effect of charged Higgs boson exchange in the leading logarithmic approximation and the heavy quark limit. In this work, we show the decay rate with QCD corrections up to the next-to-leading order (NLO). The NLO corrections are necessary to estimate theoretical uncertainties coming from short-distance calculations in the ratio of the decay rates (see below). In addition, these corrections may cause dominant uncertainties for the $`q^2`$ distribution  and the $`\tau `$ polarization .

In a two-Higgs-doublet model, the couplings of charged Higgs bosons to quarks and leptons are given by

$`_H`$ $`=`$ $`(2\sqrt{2}G_F)^{1/2}\left[X\overline{u}_LV_{KM}M_dd_R+Y\overline{u}_RM_uV_{KM}d_L+Z\overline{\nu }_LM_ll_R\right]H^+`$ (1)
$`+\text{h.c.},`$
where $`M_u`$, $`M_d`$ and $`M_l`$ are diagonal quark and lepton mass matrices, and $`V_{KM}`$ is Kobayashi-Maskawa matrix . In the MSSM, we obtain

$`X=Z=\mathrm{tan}\beta ,Y=\mathrm{cot}\beta ,`$ (2)
where $`\mathrm{tan}\beta =v_2/v_1`$ is the ratio of the vacuum expectation values of the Higgs bosons. Since the Yukawa couplings of the MSSM are the same as those of the so-called Model II of two-Higgs-doublet models , the above equations and the following results apply to the latter as well.111 It is known that SUSY loop effects in Eq.(2) are significant for large $`\mathrm{tan}\beta `$ . The dominant effect in the $`bc\tau \overline{\nu }_\tau `$ decay  comes from the SUSY-QCD correction to the bottom quark mass . Once this correction is taken into account, the Yukawa couplings of the MSSM is no longer the same as Model II of two-Higgs-doublet models. This effect cannot be ignored in order to study the MSSM Higgs sector. However, we omit it in this work because our aim is to clarify low-energy QCD corrections and uncertainties, which are universal and model-independent. The SUSY loop effects on $`\overline{B}D\tau \overline{\nu }_\tau `$ will be discussed elsewhere .

With these couplings, it turns out that the amplitude of charged Higgs exchange in $`\overline{B}D\tau \overline{\nu }_\tau `$ has a term proportional to $`m_b\mathrm{tan}^2\beta `$. Therefore, the effect of the charged Higgs boson is more significant for larger $`\mathrm{tan}\beta `$.

Formula of the decay rate is described in Sec. 2. In Sec. 3, we give hadronic form factors including next-to-leading QCD corrections. In Sec. 4, we show our numerical results. Our conclusion is given in Sec. 5.

## 2 Formula of the decay rate

Using the above Lagrangian in Eq. (1) and the standard charged current Lagrangian, we calculate the amplitudes of charged Higgs exchange and $`W`$ boson exchange in $`\overline{B}D\tau \overline{\nu }_\tau `$.

The $`W`$ boson exchange amplitude is given by

$`_s^{\lambda _\tau }(q^2,x)_W={\displaystyle \frac{G_F}{\sqrt{2}}}V_{cb}{\displaystyle \underset{\lambda _W}{}}\eta _{\lambda _W}L_{\lambda _W}^{\lambda _\tau }H_{\lambda _W}^s,`$ (3)
where $`q^2`$ is the invariant mass squared of the leptonic system, and $`x=p_Bp_\tau /m_B^2`$. The $`\tau `$ helicity and the virtual $`W`$ helicity are denoted by $`\lambda _\tau =\pm `$ and $`\lambda _W=\pm ,\mathrm{\hspace{0.17em}0},s`$, and the metric factor $`\eta _{\lambda _W}`$ is given by $`\eta _\pm =\eta _0=\eta _s=1`$. The hadronic amplitude which describes $`\overline{B}DW^{}`$ and the leptonic amplitude which describes $`W^{}\tau \overline{\nu }`$ are given by

$`H_{\lambda _W}^s(q^2)`$ $``$ $`ϵ_\mu ^{}(\lambda _W)D(p_D)|\overline{c}\gamma ^\mu b|\overline{B}(p_B),`$ (4)
$`L_{\lambda _W}^{\lambda _\tau }(q^2,x)`$ $``$ $`ϵ_\mu (\lambda _W)\tau (p_\tau ,\lambda _\tau )\overline{\nu }_\tau (p_\nu )|\overline{\tau }\gamma ^\mu (1\gamma _5)\nu _\tau |0,`$ (5)
where $`ϵ_\mu (\lambda _W)`$ is the polarization vector of the virtual $`W`$ boson.

The charged Higgs exchange amplitude is given by

$`_s^{\lambda _\tau }(q^2,x)_H={\displaystyle \frac{G_F}{\sqrt{2}}}V_{cb}L^{\lambda _\tau }\left[XZ^{}{\displaystyle \frac{m_bm_\tau }{M_H^2}}H_R^s+YZ^{}{\displaystyle \frac{m_cm_\tau }{M_H^2}}H_L^s\right].`$ (6)
Here, the hadronic and leptonic amplitudes are defined by

$`H_{R,L}^s(q^2)`$ $``$ $`D(p_D)|\overline{c}(1\pm \gamma _5)b|\overline{B}(p_B),`$ (7)
$`L^{\lambda _\tau }(q^2,x)`$ $``$ $`\tau (p_\tau ,\lambda _\tau )\overline{\nu }_\tau (p_\nu )|\overline{\tau }(1\gamma _5)\nu _\tau |0.`$ (8)
This leptonic amplitude is related to the $`W`$ exchange amplitude as

$`L^{\lambda _\tau }={\displaystyle \frac{\sqrt{q^2}}{m_\tau }}L_s^{\lambda _\tau }.`$ (9)
Details of the hadronic amplitudes for the $`W`$ exchange and the charged Higgs exchange are discussed in the next section.

Using the amplitudes of Eqs. (3) and (6), the differential decay rate is given by

$`{\displaystyle \frac{d\mathrm{\Gamma }}{dq^2}}`$ $`=`$ $`{\displaystyle \frac{G_F^2|V_{cb}|^2\overline{v}^4\sqrt{Q_+Q_{}}}{128\pi ^3m_B^3}}[({\displaystyle \frac{2}{3}}q^2+{\displaystyle \frac{1}{3}}m_\tau ^2)(H_0^s)^2`$ (10)
$`+m_\tau ^2((m_b\mathrm{tan}^2\beta +m_c){\displaystyle \frac{\sqrt{q^2}}{M_H^2}}H_R^sH_s^s)^2],`$
where $`Q_\pm =(m_B\pm m_D)^2q^2`$ and $`\overline{v}=\sqrt{1m_\tau ^2/q^2}`$. Note that if $`\mathrm{tan}\beta 1`$, in which we are interested, this decay rate is practically a function of $`\mathrm{tan}\beta /M_H`$ because the second term in the coefficient of $`H_R^s`$ is negligible for $`m_b\mathrm{tan}^2\beta m_c`$.

## 3 Hadronic form factors including QCD corrections

Here, we evaluate the hadronic amplitudes in Eq. (4) and Eq. (7) in order to obtain the decay rate numerically. These amplitudes are given in terms of hadronic form factors:

$`D(p_D)|\overline{c}\gamma ^\mu b|\overline{B}(p_B)`$ $`=`$ $`\sqrt{m_Bm_D}\left[h_+(w)(v+v^{})^\mu +h_{}(w)(vv^{})^\mu \right],`$ (11)
$`D(p_D)|\overline{c}b|\overline{B}(p_B)`$ $`=`$ $`\sqrt{m_Bm_D}(1+w)h_s(w),`$ (12)
where $`v=p_B/m_B`$, $`v^{}=p_D/m_D`$ and $`wvv^{}=(m_B^2+m_D^2q^2)/2m_Bm_D`$ .

In the heavy quark limit and in the leading logarithmic approximation, these form factors $`h_\pm (w)`$ and $`h_s(w)`$ are given as

$`h_+(w)=h_s(w)=\xi (w),h_{}(w)=0,`$ (13)
where $`\xi (w)`$ is the universal form factor .

Now, we consider QCD corrections beyond LLA and calculate these form factors up to the next-to-leading order. Then, these form factors are given as

$`h_+(w)`$ $`=`$ $`\left\{\widehat{C}_1(w)\left({\displaystyle \frac{w+1}{2}}\right)\left(\widehat{C}_2(w)+\widehat{C}_3(w)\right)\right\}\xi (w),`$ (14)
$`h_{}(w)`$ $`=`$ $`\left\{{\displaystyle \frac{w+1}{2}}\left(\widehat{C}_2(w)\widehat{C}_3(w)\right)\right\}\xi (w),`$ (15)
$`h_s(w)`$ $`=`$ $`\widehat{C}_s(w)\xi (w).`$ (16)
Explicit formula of coefficients for $`W`$ exchange, $`\widehat{C}_i(i=13)`$, are given by Neubert . The coefficient for charged Higgs exchange, $`\widehat{C}_s`$, is given as

$`\widehat{C}_s(w)=A(w)\overline{C}_s(w),`$ (17)
where

$`A(w)`$ $`=`$ $`\left({\displaystyle \frac{\alpha _s(m_c)}{\alpha _s(m_b)}}\right)^{\frac{6}{25}}[\alpha _s(m_c)]^{a_L(w)},`$ (18)
$`\overline{C}_s(w)`$ $`=`$ $`1+{\displaystyle \frac{\alpha _s(m_b)\alpha _s(m_c)}{\pi }}(\stackrel{~}{Z}+2)+{\displaystyle \frac{\alpha _s(m_c)}{\pi }}\left[Z(w)+2+{\displaystyle \frac{3}{2}}\left(f(w)r(w)\right)\right]`$ (19)
$`+{\displaystyle \frac{2\alpha _s(\overline{m})}{3\pi }}g_s(z,w).`$
$`\overline{m}`$ is some average mass of $`m_b`$ and $`m_c`$ and other functions and constants in Eq. (18) and Eq. (19) are given in Ref. and Appendix. We have used the $`\overline{\mathrm{MS}}`$ scheme in our calculations.

The form of $`\xi (w)`$ is strongly constrained by the dispersion relations as

$`\xi (w)18\rho _1^2z+(51.\rho _1^210.)z^2(252.\rho _1^284.)z^3,`$ (20)
where $`z=(\sqrt{w+1}\sqrt{2})/(\sqrt{w+1}+\sqrt{2})`$. We obtain the slope parameter $`\rho _1^2`$ as

$`\rho _1^2=1.33\pm 0.22,`$ (21)
from the experimental data of $`BD^{}e\overline{\nu }_e`$ .

## 4 Numerical results

We consider the ratio of decay rates,

$`B={\displaystyle \frac{\mathrm{\Gamma }(\overline{B}D\tau \overline{\nu }_\tau )}{\mathrm{\Gamma }(\overline{B}D\mu \overline{\nu }_\mu )_{SM}}},`$ (22)
where the denominator is the decay rate of $`\overline{B}D\mu \overline{\nu }_\mu `$ in the SM. Uncertainties due to the form factors and other parameters tend to reduce or vanish by taking the ratio.

Fig. 1(a) is the plot of our prediction of the ratio in Eq. (22) as a function of $`R`$, which is defined by $`Rm_W\mathrm{tan}\beta /m_H`$. Here, we do not show the error in the slope parameter $`\rho _1^2`$ and we take $`\rho _1^2=1.33`$. The dashed lines show the MSSM and SM predictions without QCD corrections. The lines with narrow shaded regions show the predictions including QCD corrections with $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ being varied between 0.15 GeV and 0.25 GeV.

In Fig. 1(b) we show the ratio of $`B`$ in the MSSM with and without QCD corrections, $`B`$(with QCD corrections)/$`B`$(without QCD corrections), as a function of $`R`$. The solid line is the ratio with $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}=0.25`$ GeV and the dashed line is that with $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}=0.15`$ GeV.

From Fig. 1, we expect that theoretical uncertainties from higher order QCD corrections are at most a few percents in the ratio of decay rates and, as seen later, the theoretical uncertainties from QCD corrections are much smaller than those from the error of $`\rho _1^2`$.

Fig. 2(a) shows our prediction of the ratio $`B`$ with QCD corrections as a function of $`R`$. Here, we take the error in the slope parameter into account and use $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}=0.25`$ GeV. The shaded regions show the MSSM and SM predictions with the error in the slope parameter $`\rho _1^2`$ in Eq. (21). As mentioned before, from Fig. 2(a) and Fig. 1, we see that the theoretical uncertainty from the error in the slope parameter is dominant over that from QCD corrections. As seen in Fig. 1(a), when $`R`$ is about 35, the ratio $`B`$ in the MSSM is the same as the one in the SM.

In Fig. 2(b), we also show the ratio,

$`\stackrel{~}{B}={\displaystyle \frac{\mathrm{\Gamma }(\overline{B}D\tau \overline{\nu }_\tau )}{\stackrel{~}{\mathrm{\Gamma }}(\overline{B}D\mu \overline{\nu }_\mu )_{SM}}},`$ (23)
with QCD corrections, similar as Fig. 2(a), but its denominator is $`\stackrel{~}{\mathrm{\Gamma }}(\overline{B}D\mu \overline{\nu }_\mu )_{SM}`$, which is integrated over the same $`q^2`$ region as the $`\tau `$ mode, i.e., $`m_\tau ^2q^2(m_Bm_D)^2`$. From Fig. 2(b), we observe that the ratio $`\stackrel{~}{B}`$ has less theoretical uncertainty and we expect a better sensitivity than the ratio $`B`$ in Fig. 2(a).

Now, we consider a decay distribution defined by

$`d(w)=(w^21){\displaystyle \frac{\mathrm{\Gamma }(\overline{B}D\tau \overline{\nu }_\tau )/dw}{\mathrm{\Gamma }(\overline{B}D\mu \overline{\nu }_\mu )_{SM}/dw}},`$ (24)
where $`w=(m_B^2+m_D^2q^2)/2m_Bm_D`$.

In Fig. 3, we show $`d(w)`$ in the SM and the MSSM with different values of $`R`$. The dashed lines show the MSSM and SM predictions without QCD corrections. The lines with shaded regions show the predictions with QCD corrections and the shaded regions are given by using $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}=0.150.25`$ GeV. The theoretical uncertainty from the error in the slope parameter is canceled out in this quantity. Thus, QCD corrections become dominant uncertainties in the theoretical calculation. From Fig. 2, the ratio $`B`$ in the MSSM becomes the same as the one in the SM when $`R35`$. But, in Fig. 3, we find that the behavior of $`d(w)`$ in the MSSM with $`R=35`$ is considerably different from that in the SM. Therefore, we can distinguish the SM from the MSSM by investigating $`w`$ distribution even if $`R35`$.

## 5 Conclusion

As seen in our numerical results, the branching ratio of $`\overline{B}D\tau \overline{\nu }_\tau `$ is a sensitive probe of the MSSM-like Higgs sector. So, if $`\overline{B}D\tau \overline{\nu }_\tau `$ is observed at a B factory experiment, a significant region of the parameter space of the MSSM Higgs sector will be proved. This is complementary at the Higgs search at LHC .

In the branching ratio, the theoretical uncertainty from QCD correction is much smaller than that from the error in the slope parameter $`\rho _1^2`$. However, in $`w`$ distribution, the theoretical uncertainty from the error in the slope parameter $`\rho _1^2`$ is canceled out and, therefore, it is important to consider QCD corrections.

In this work, we have not taken $`1/m`$ corrections into account. They may be as significant as QCD corrections in the $`w`$ distribution. However we expect that their effects are smaller than the uncertainty from the error of $`\rho _1^2`$ in the branching ratio .

Finally, the $`\tau `$ polarization in $`\overline{B}D\tau \overline{\nu }_\tau `$ is also expected to be a good probe of charged Higgs boson. The theoretical uncertainty from the error in the slope parameter becomes very small in the $`\tau `$ polarization . QCD and $`1/m`$ corrections to this quantity will be addressed elsewhere.

## Appendix

Functions in Eq. (18) and Eq. (19) are given by

$`r(w)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{w^21}}}\mathrm{ln}(w+\sqrt{w^21}),`$ (A.1)
$`f(w)`$ $`=`$ $`wr(w)2{\displaystyle \frac{w}{2\sqrt{w^21}}}\left[L_2(1w_{}^2)L_2(1w_+^2)\right],`$ (A.2)
$`g_s(z,w)`$ $`=`$ $`{\displaystyle \frac{w}{\sqrt{w^21}}}\left[L_2(1zw_{})L_2(1zw_+)\right]`$ (A.3)
$`{\displaystyle \frac{z}{12wz+z^2}}\left[(w^21)r(w)+(wz)\mathrm{ln}z\right],`$
$`a_L(w)`$ $`=`$ $`{\displaystyle \frac{8}{27}}\left[wr(w)1\right],`$ (A.4)
$`Z(w)`$ $`=`$ $`{\displaystyle \frac{4}{9}}\left[{\displaystyle \frac{25}{54}}+{\displaystyle \frac{\pi ^2}{12}}+{\displaystyle \frac{5}{9}}+f(w)\right]\left[wr(w)1\right]{\displaystyle \frac{8}{9}}I(w),`$ (A.5)
$`\stackrel{~}{Z}`$ $`=`$ $`{\displaystyle \frac{7}{225}}\pi ^2{\displaystyle \frac{9403}{7500}},`$ (A.6)
where $`z=m_c/m_b`$, $`w_\pm =w\pm \sqrt{w^21}`$,

$`L_2(x)={\displaystyle _0^x}𝑑t{\displaystyle \frac{\mathrm{ln}(1t)}{t}},`$ (A.7)
is the dilogarithm and,

$`I(w)`$ $`=`$ $`{\displaystyle _0^\phi }𝑑\psi \left[\psi \mathrm{coth}\psi 1\right]`$ (A.8)
$`\times \left\{\psi \mathrm{coth}^2\phi +{\displaystyle \frac{\mathrm{sinh}\phi \mathrm{cosh}\phi }{\mathrm{sinh}^2\phi \mathrm{sinh}^2\psi }}\mathrm{ln}{\displaystyle \frac{\mathrm{sinh}\phi }{\mathrm{sinh}\psi }}\right\},`$
in terms of the hyperbolic angle $`\phi `$ defined by $`w=\mathrm{cosh}\phi `$.


text:
# Theorem 1.1

Approximation to real numbers by cubic algebraic integers I

by Damien ROY

2000 Mathematics Subject Classification: Primary 11J04; Secondary 11J13, 11J82 Work partly supported by NSERC and CICMA.
1. Introduction

The study of approximation to a real number by algebraic numbers of bounded degree started with a paper of E. Wirsing  in 1961. Motivated by this, H. Davenport and W. M. Schmidt considered in  the analogous inhomogeneous problem of approximation to a real number by algebraic integers of bounded degree. They proved a result that is optimal for degree $`2`$ and a general result which is valid for any degree. We shall be concerned here with the case of degree $`3`$, the first case for which the optimal exponent of approximation is not known.

To state their results, we define the height $`H(P)`$ of a polynomial $`P𝐑[T]`$ to be the maximum of the absolute values of its coefficients, and define the height $`H(\alpha )`$ of an algebraic number $`\alpha `$ to be the height of its irreducible polynomial over $`𝐙`$. We also denote by $`\gamma =(1+\sqrt{5})/2`$ the golden ratio.

Davenport and Schmidt showed that, for any real number $`\xi `$ which is neither rational nor quadratic irrational, there exist a positive constant $`c`$ depending only on $`\xi `$ and arbitrarily large real numbers $`X`$ such that the inequalities

$$|x_0|X,|x_0\xi x_1|cX^{1/\gamma },|x_0\xi ^2x_2|cX^{1/\gamma }$$
$`(1.1)`$
have no solutions in integers $`x_0`$, $`x_1`$, $`x_2`$, not all zero (see Theorem 1a of ). By an argument of duality based on geometry of numbers, they deduced from this that, for such a real number $`\xi `$, there exist another constant $`c>0`$ and infinitely many algebraic integers $`\alpha `$ of degree at most $`3`$ which satisfy

$$0<|\xi \alpha |cH(\alpha )^{\gamma ^2}$$
$`(1.2)`$
(see Theorem 1 of ).

At first, it would be natural to expect that the first statement holds with any exponent $`>1/2`$ in place of $`1/\gamma 0.618`$ in the inequalities (1.1), since this ensures that the volume of the convex body defined by these inequalities tends to zero as $`X`$ tends to infinity. However, this is not the case, and we will show in fact that the conditions (1.1) in this statement are optimal (see also  for an announcement).

###### Theorem 1.1

There exists a real number $`\xi `$ which is neither rational nor quadratic irrational and which has the property that, for a suitable constant $`c>0`$, the inequalities (1.1) have a non-zero solution $`(x_0,x_1,x_2)𝐙^3`$ for any real number $`X1`$. Any such number is transcendental over $`𝐐`$ and the set of these real numbers is countable.

Let us call extremal any real number that has the property stated in Theorem 1.1. For such a number $`\xi `$, solutions to (1.1) provide simultaneous approximations of $`\xi `$ and $`\xi ^2`$ by rational numbers $`x_1/x_0`$ and $`x_2/x_0`$ with the same denominator (assuming $`x_00`$). The fact that, for any sufficiently large $`X`$, these are sharper than expected by an application of Dirichlet box principle (see for example Theorem 1A in Chapter II of ) reminds of Cassels’ counter-example in questions of algebraic independence (see Theorem XIV in Chapter V of ). Another situation, even closer to that of Cassels, where these numbers satisfy estimates that are better than expected from the box principle for all sufficiently large value of the parameter $`X`$ is the following (compare with Theorem 1B in Chapter II of ).

###### Theorem 1.2

Let $`\xi `$ be an extremal real number. There exists a constant $`c>0`$ depending only on $`\xi `$ such that for any real number $`X1`$ the inequalities

$$|x_0|X,|x_1|X,|x_0\xi ^2+x_1\xi +x_2|cX^{\gamma ^2}$$
$`(1.3)`$
have a non-zero solution $`(x_0,x_1,x_2)𝐙^3`$.

Indeed, since $`\gamma ^22.618>2`$, the volume of the convex body of $`𝐑^3`$ defined by (1.3) tends to zero as $`X`$ tends to infinity, a situation where we do not expect in principle to find a non-zero solution for each sufficiently large $`X`$.

Examples of extremal real numbers are the numbers $`\xi _{a,b}`$ whose continued fraction expansion $`\xi _{a,b}=[0,a,b,a,a,b,a,b,\mathrm{}]`$ is given by the Fibonacci word on two distinct positive integers $`a`$ and $`b`$ (the infinite word $`abaabab\mathrm{}`$ starting with $`a`$ which is the fixed point of the substitution sending $`a`$ to $`ab`$ and $`b`$ to $`a`$). We don’t know if, like these numbers $`\xi _{a,b}`$, all extremal real numbers have bounded partial quotients but they satisfy the following weaker measure of approximation by rational numbers.

###### Theorem 1.3

Let $`\xi `$ be an extremal real number. Then, there exist constants $`c>0`$ and $`t>0`$ such that, for any rational number $`\alpha 𝐐`$, we have

$$|\xi \alpha |cH(\alpha )^2(1+\mathrm{log}H(\alpha ))^t.$$

The above mentioned Fibonacci continued fractions $`\xi _{a,b}`$ are an example of the more general Sturmian continued fractions studied by J.-P. Allouche, J. L. Davison, M. Queffélec and L. Q. Zamboni in . These authors showed that Sturmian continued fractions are transcendental by showing that they admit very good approximations by quadratic real numbers. In the present situation, we get an optimal measure of approximation by such numbers.

###### Theorem 1.4

Let $`\xi `$ be an extremal real number. There is a value of $`c>0`$ such that the inequality

$$|\xi \alpha |cH(\alpha )^{2\gamma ^2}$$

has infinitely many solutions in algebraic numbers $`\alpha 𝐂`$ of degree $`2`$ and another value of $`c>0`$ such that the same inequality has no solution in algebraic numbers $`\alpha 𝐂`$ of degree at most $`2`$.

It is generally believed that, for any real number $`\xi `$ that is not algebraic over $`𝐐`$ of degree at most $`3`$ and for any $`ϵ>0`$, there should exist infinitely many algebraic integers $`\alpha `$ of degree at most $`3`$ which satisfy $`|\xi \alpha |H(\alpha )^{3+ϵ}`$ (see page 259 of ). A recent result of Y. Bugeaud and O. Teulié shows that this is true of almost all real numbers $`\xi `$ in the sense of Lebesgue measure (see Corollary 1 to Theorem 6 of ). By Schmidt’s subspace theorem, this is also true of all algebraic numbers of degree $`>3`$. However, we will present a criterion which suggests that this is probably not true for extremal real numbers and that, for these numbers, the correct exponent of approximation is probably $`\gamma ^2`$ instead of $`3`$ ($``$) Since this paper was written the author found a class of extremal real numbers for which this exponent of approximation is indeed $`\gamma ^2`$; see Approximation to real numbers by cubic algebraic integers II, Annals of Math. (to appear), arXiv:math.NT/0210182.. As an application of this criterion, we will derive the following measure of approximation.

###### Theorem 1.5

Let $`\xi `$ be an extremal real number. There exists a constant $`c>0`$ such that, for any algebraic integer $`\alpha 𝐂`$ of degree $`3`$, we have

$$|\xi \alpha |cH(\alpha )^{\gamma ^21}.$$

This does not disprove the natural conjecture stated above, but it sheds doubts on it. By extension, it also sheds doubts on the conjectures of Wirsing and Schmidt concerning approximation to a real number by algebraic numbers of bounded degree, although Davenport and Schmidt have established these conjectures in the case of approximation by quadratic irrationals in .

This paper is organized as follows. Section 5 presents a criterion for a real number to be extremal. Its proof is based on preliminary considerations from Sections 2, 3 and 4. Section 6 provides examples of extremal real numbers which include the Fibonacci continued fractions. Sections 7, 8 and 9 are devoted to measures of approximations of extremal real numbers by rational numbers, by quadratic real numbers and by algebraic integers of degree at most 3 respectively. Section 9 also provides a criterion for approximation by algebraic integers of degree $`3`$, and Section 10 proves a partial converse to this criterion.

As suggested by the referee, the reader may, at a first reading, skip Section 2 and part (iii) of Lemma 3.1 and go directly from Theorem 5.1 to Section 7. He may also want to devise an independent proof of Proposition 6.3 showing that Fibonacci continued fractions are extremal real numbers, or look at the short note  for such a proof. However, in a more linear reading, the relevance of the algebraic tool developed in Section 2 will appear clearly in Corollary 5.2, a result which the author sees as a key for further study of extremal real numbers and a motivation for the constructions of Section 6.

Notation. Given a fixed real number $`\xi `$, we often write estimates of the form $`AB`$ or $`BA`$, where $`A`$ and $`B`$ are non-negative variable real numbers, to mean $`AcB`$ for some constant $`c>0`$ depending only on $`\xi `$. We also write $`AB`$ to mean $`ABA`$.

2. A particular construction

Let $`A`$ denote a commutative ring with unit (in practice, we will take $`A=𝐙`$). Identifying any point $`𝐱=(x_0,x_1,x_2)A^3`$ with the symmetric matrix

$$𝐱=\left(\begin{array}{cc}x_0& x_1\\ x_1& x_2\end{array}\right),$$

we can talk of its determinant

$$det(𝐱)=x_0x_2x_1^2.$$

Put

$$J=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).$$

Then, given any three points $`𝐱,𝐲,𝐳A^3`$, we find, with the natural notation for their coordinates,

$$\begin{array}{cc}\hfill 𝐱J𝐳J𝐲& =\left(\begin{array}{cc}x_0& x_1\\ x_1& x_2\end{array}\right)\left(\begin{array}{cc}z_2& z_1\\ z_1& z_0\end{array}\right)\left(\begin{array}{cc}y_0& y_1\\ y_1& y_2\end{array}\right)\hfill \\ & \\ & =\left(\begin{array}{cc}x_0& x_1\\ x_1& x_2\end{array}\right)\left(\begin{array}{cc}\left|\begin{array}{cc}y_0& y_1\\ z_1& z_2\end{array}\right|& \left|\begin{array}{cc}y_1& y_2\\ z_1& z_2\end{array}\right|\\ \\ \\ \left|\begin{array}{cc}y_0& y_1\\ z_0& z_1\end{array}\right|& \left|\begin{array}{cc}y_1& y_2\\ z_0& z_1\end{array}\right|\end{array}\right)\hfill \\ & \\ & =\left(\begin{array}{cc}\left|\begin{array}{cc}x_0& x_1\\ \\ \left|\begin{array}{cc}y_0& y_1\\ z_0& z_1\end{array}\right|& \left|\begin{array}{cc}y_0& y_1\\ z_1& z_2\end{array}\right|\end{array}\right|& \left|\begin{array}{cc}x_0& x_1\\ \\ \left|\begin{array}{cc}y_1& y_2\\ z_0& z_1\end{array}\right|& \left|\begin{array}{cc}y_1& y_2\\ z_1& z_2\end{array}\right|\end{array}\right|\\ \\ \\ \left|\begin{array}{cc}x_1& x_2\\ \\ \left|\begin{array}{cc}y_0& y_1\\ z_0& z_1\end{array}\right|& \left|\begin{array}{cc}y_0& y_1\\ z_1& z_2\end{array}\right|\end{array}\right|& \left|\begin{array}{cc}x_1& x_2\\ \\ \left|\begin{array}{cc}y_1& y_2\\ z_0& z_1\end{array}\right|& \left|\begin{array}{cc}y_1& y_2\\ z_1& z_2\end{array}\right|\end{array}\right|\end{array}\right).\hfill \end{array}$$

Moreover, if, in this matrix, we subtract from the element of row 1 and column 2 the element of row 2 and column 1, we find a polynomial in $`𝐱`$, $`𝐲`$ and $`𝐳`$ which happens to be

$$\mathrm{trace}(J𝐱J𝐳J𝐲)=det(𝐱,𝐲,𝐳).$$
$`(2.1)`$
So, the product $`𝐱J𝐳J𝐲`$ is symmetric if and only if $`det(𝐱,𝐲,𝐳)=0`$. In this case, we denote by $`[𝐱,𝐲,𝐳]`$ the corresponding element of $`A^3`$:

$$[𝐱,𝐲,𝐳]=𝐱J𝐳J𝐲.$$
$`(2.2)`$
###### Lemma 2.1

For any $`𝐱,𝐲,𝐳A^3`$ with $`det(𝐱,𝐲,𝐳)=0`$ and any $`𝐰A^3`$, we have:

(i) $`det[𝐱,𝐲,𝐳]=det(𝐱)det(𝐲)det(𝐳)`$,

(ii) $`det(𝐰,𝐲,[𝐱,𝐲,𝐳])=det(𝐲)det(𝐰,𝐳,𝐱)`$,

(iii) $`det(𝐱,𝐲,[𝐱,𝐲,𝐳])=0`$ and $`[𝐱,𝐲,[𝐱,𝐲,𝐳]]=det(𝐱)det(𝐲)𝐳`$.

Proof. The formula (i) follows from (2.2) since $`det(J)=1`$. To prove (ii) and (iii), we note that, for any $`𝐰A^3`$, we have

$$𝐰J𝐰J=J𝐰J𝐰=det(𝐰)I$$

where $`I`$ denotes the identity matrix. Combining this observation with (2.1), we find

$$\begin{array}{cc}\hfill det(𝐰,𝐲,[𝐱,𝐲,𝐳])& =\mathrm{trace}(J𝐰J(𝐱J𝐳J𝐲)J𝐲)\hfill \\ & =det(𝐲)\mathrm{trace}(J𝐰J𝐱J𝐳)\hfill \\ & =det(𝐲)det(𝐰,𝐳,𝐱).\hfill \end{array}$$

This proves (ii) and the first half of (iii). For the second half, we compute

$$𝐱J(𝐱J𝐳J𝐲)J𝐲=(𝐱J𝐱J)𝐳(J𝐲J𝐲)=det(𝐱)det(𝐲)𝐳.$$

3. Some estimates

Let $`\xi `$ be any real number. For each integral point $`𝐱=(x_0,x_1,x_2)𝐙^3`$, we define

$$𝐱=\mathrm{max}\{|x_0|,|x_1|,|x_2|\}\text{and}L(𝐱)=\mathrm{max}\{|x_0\xi x_1|,|x_0\xi ^2x_2|\}.$$

For a square matrix $`A`$, we also denote by $`A`$ the absolute value of its determinant. In the following lemma, we recall estimates for determinants that appear in the study of Davenport and Schmidt in Section 3 of , and we provide estimates relative to the construction of the previous paragraph.

###### Lemma 3.1

Let $`𝐱,𝐲,𝐳𝐙^3`$.

(i) For any choice of integers $`r,s,t,u\{0,1,2\}`$ with $`sr=ut`$, we have

$$\begin{array}{cc}x_r& x_s\\ y_t& y_u\end{array}𝐱L(𝐲)+𝐲L(𝐱).$$

(ii) We have $`|det(𝐱,𝐲,𝐳)|𝐱L(𝐲)L(𝐳)+𝐲L(𝐱)L(𝐳)+𝐳L(𝐱)L(𝐲).`$

(iii) Put $`𝐰=[𝐱,𝐱,𝐲]`$. Then, we have

$$𝐰𝐱^2L(𝐲)+𝐲L(𝐱)^2\text{and}L(𝐰)\left(𝐱L(𝐲)+𝐲L(𝐱)\right)L(𝐱).$$

Proof. All these estimates are based on the multilinearity of the determinant. We simply prove (iii) since (i) and (ii) follow from the computations in the proofs of Lemmas 3 and 4 of . Write $`𝐰=(w_0,w_1,w_2)`$. Since

$$w_0=\left|\begin{array}{cc}x_0& x_1\\ \\ \left|\begin{array}{cc}x_0& x_1\\ y_0& y_1\end{array}\right|& \left|\begin{array}{cc}x_0& x_1\\ y_1& y_2\end{array}\right|\end{array}\right|=\left|\begin{array}{cc}x_0& x_1x_0\xi \\ \\ \left|\begin{array}{cc}x_0& x_1\\ y_0& y_1\end{array}\right|& \left|\begin{array}{cc}x_0& x_1\\ y_1y_0\xi & y_2y_1\xi \end{array}\right|\end{array}\right|,$$

we obtain, using part (i) and noting that $`L(𝐱)𝐱`$,

$$|w_0|𝐱\begin{array}{cc}x_0& x_1\\ y_1y_0\xi & y_2y_1\xi \end{array}+L(𝐱)\begin{array}{cc}x_0& x_1\\ y_0& y_1\end{array}𝐱^2L(𝐲)+𝐲L(𝐱)^2.$$

Similarly, we find, for $`j=0,1`$,

$$|w_{j+1}w_j\xi |=\begin{array}{cc}x_1x_0\xi & x_2x_1\xi \\ \\ \left|\begin{array}{cc}x_j& x_{j+1}\\ y_0& y_1\end{array}\right|& \left|\begin{array}{cc}x_j& x_{j+1}\\ y_1& y_2\end{array}\right|\end{array}L(𝐱)(𝐱L(𝐲)+𝐲L(𝐱))$$

and the conclusion follows.

4. The sequence of minimal points

Let $`\xi `$ be a real number which is neither rational nor quadratic over $`𝐐`$. Then, the function $`L:𝐙^3𝐑`$ defined in Section 3 has the property that, for $`𝐱,𝐲𝐙^3`$ with $`x_0,y_00`$, we have $`L(𝐱)=L(𝐲)`$ if and only if $`𝐱=\pm 𝐲`$. Accordingly, for any real number $`X1`$, the finite set of points $`𝐱=(x_0,x_1,x_2)`$ of $`𝐙^3`$ satisfying

$$1x_0X\text{and}L(𝐱)=\mathrm{max}\{|x_0\xi x_1|,|x_0\xi ^2x_2|\}1$$

contains exactly one point for which $`L(𝐱)`$ is minimal. Following Davenport and Schmidt, we call it the minimal point corresponding to $`X`$.

Clearly, if $`𝐱`$ is a minimal point for some real number $`X1`$, then it is a minimal point for $`x_0`$. So, we may order the minimal points according to their first coordinates. Note also that the coordinates of a minimal point are relatively prime as a set. Therefore, distinct minimal points are linearly independent over $`𝐐`$ as elements of $`𝐐^3`$. In particular, any two consecutive minimal points are linearly independent over $`𝐐`$.

To state the next lemma, we define the height of a $`2\times 3`$ matrix $`A`$, denoted $`A`$, to be the maximal absolute value of its minors of order $`2`$, and we define the height of a $`2`$-dimensional subspace $`V`$ of $`𝐐^3`$, denoted $`H(V)`$, to be the height of any $`2\times 3`$ matrix whose rows form a basis of $`V𝐙^3`$.

###### Lemma 4.1

Let $`𝐱`$ be a minimal point, let $`𝐲`$ be the next minimal point and let $`V=𝐱,𝐲_𝐐`$ be the subspace of $`𝐐^3`$ generated by $`𝐱`$ and $`𝐲`$. Then, $`\{𝐱,𝐲\}`$ is a basis of $`V𝐙^3`$ over $`𝐙`$ and we have $`H(V)𝐲L(𝐱)`$.

Proof. Write $`𝐱=(x_0,x_1,x_2)`$ and $`𝐲=(y_0,y_1,y_2)`$. Arguing like Davenport and Schmidt in their proof of Lemma 2 of , we observe that, if $`\{𝐱,𝐲\}`$ were not a basis of $`V𝐙^3`$, there would exist a point $`𝐳=(z_0,z_1,z_2)`$ of $`𝐙^3`$ with $`z_0>0`$, which could be written in the form $`𝐳=r𝐱+s𝐲`$ for rational numbers $`r`$ and $`s`$ with $`\mathrm{max}\{|r|,|s|\}1/2`$. We would then get

$$z_0|r|x_0+|s|y_0<y_0\text{and}L(𝐳)|r|L(𝐱)+|s|L(𝐲)<L(𝐱),$$

in contradiction with the fact that $`𝐲`$ is the next minimal point after $`𝐱`$. Thus, $`\{𝐱,𝐲\}`$ is a basis of $`V`$ and, using Lemma 3.1 (i), we find

$$H(V)=\begin{array}{ccc}x_0& x_1& x_2\\ y_0& y_1& y_2\end{array}𝐲L(𝐱)+𝐱L(𝐲)𝐲L(𝐱).$$

To derive lower bounds for $`H(V)`$, put $`u_i=x_ix_0\xi ^i`$ and $`v_i=y_iy_0\xi ^i`$ for $`i=1,2`$, and choose an index $`j`$ for which $`|u_j|=L(𝐱)`$. We find

$$H(V)\begin{array}{cc}x_0& x_j\\ y_0& y_j\end{array}=|x_0v_jy_0u_j|y_0|u_j|x_0|v_j|(y_0x_0)L(𝐱).$$

Since the point $`𝐳=𝐲𝐱`$ has its first coordinate $`z_0=y_0x_0`$ satisfying $`1z_0<y_0`$, and since $`𝐲`$ is the next minimal point after $`𝐱`$, there must exist an index $`i`$ for which $`|z_iz_0\xi ^i|=|v_iu_i|>L(𝐱)`$. Thus, we also find

$$H(V)|x_0v_iy_0u_i||x_0(v_iu_i)(y_0x_0)u_i|x_0L(𝐱)(y_0x_0)L(𝐱).$$

Combining the two preceding displayed inequalities, we find

$$3H(V)x_0L(𝐱)+(y_0x_0)L(𝐱)y_0L(𝐱),$$

and therefore $`H(V)𝐲L(𝐱)`$.

5. The set of extremal numbers

Recall that, in the introduction, we defined a real number $`\xi `$ to be extremal if it is not rational nor quadratic irrational and if there exists a constant $`c>0`$ such that the inequalities

$$1|x_0|X,|x_0\xi x_1|cX^{1/\gamma },|x_0\xi ^2x_2|cX^{1/\gamma },$$
$`(5.1)`$
have a solution in integers $`x_0,x_1,x_2`$ for any real number $`X1`$. By virtue of the subspace theorem of Schmidt, such a real number is transcendental over $`𝐐`$ (see for example Theorem 1B in Chapter VI of ).

In this section, we characterize the extremal real numbers by a stronger approximation property and show that the corresponding approximation triples satisfy a certain recurrence relation involving the operation $`[𝐱,𝐲,𝐳]`$ defined in Section 2. We conclude that the set of extremal numbers is at most countable.

###### Theorem 5.1

A real number $`\xi `$ is extremal if and only if there exists an increasing sequence of positive integers $`(Y_k)_{k1}`$ and a sequence of points $`(𝐲_k)_{k1}`$ of $`𝐙^3`$ such that, for all $`k1`$, we have

$$Y_{k+1}Y_k^\gamma ,𝐲_kY_k,L(𝐲_k)Y_k^1,$$

$$1|det(𝐲_k)|1\text{and}1|det(𝐲_k,𝐲_{k+1},𝐲_{k+2})|1,$$

where the norm $``$ and the map $`L=L_\xi `$ are as defined in Section 3.

Proof. Assume first that there exist such sequences $`(𝐲_k)_{k1}`$ and $`(Y_k)_{k1}`$ for a given real number $`\xi `$. Then, $`\xi `$ is not rational nor quadratic irrational because otherwise there would exist integers $`p,q,r`$ not all zero such that $`p+q\xi +r\xi ^2=0`$ and consequently, for any point $`𝐲_k=(y_{k,0},y_{k,1},y_{k,2})`$, we would have

$$|py_{k,0}+qy_{k,1}+ry_{k,2}|=|q(y_{k,1}y_{k,0}\xi )+r(y_{k,2}y_{k,0}\xi ^2)|Y_k^1.$$

Since the left hand side of this inequality is an integer, this integer would therefore be zero for all sufficiently large $`k`$, in contradiction with the fact that the determinant of any three consecutive points of the sequence $`(𝐲_k)_{k1}`$ is non-zero. Moreover, for any sufficiently large real number $`X`$, there exists an index $`k1`$ such that $`Y_kX<Y_{k+1}`$ and then the point $`𝐲_k`$ satisfies

$$𝐲_kY_kX\text{and}L(𝐲_k)Y_k^1Y_{k+1}^{1/\gamma }X^{1/\gamma }.$$

Therefore, $`\xi `$ is extremal.

Conversely, assume that $`\xi `$ is an extremal real number and consider the sequence $`(𝐱_i)_{i1}`$ of minimal points of $`𝐙^3`$ associated to $`\xi `$, ordered according to their first coordinate (see Section 4 above). For each $`i1`$, denote by $`X_i`$ the first coordinate of $`𝐱_i`$ and put $`L_i=L(𝐱_i)`$. Choose also $`c>0`$ such that (5.1) has an integral solution for each $`X1`$. Then, by the definition of minimal points, we have

$$L_icX_{i+1}^{1/\gamma },(i1)$$
$`(5.2)`$
(see formula (30) in ). By Lemma 2 of , there is an index $`i_02`$ such that $`det(𝐱_i)0`$ for all $`ii_0`$. Moreover, Lemma 5 of  shows that there are infinitely many indices $`i>i_0`$ such that $`𝐱_{i1}`$, $`𝐱_i`$ and $`𝐱_{i+1}`$ are linearly independent over $`𝐐`$. Let $`I`$ denote the set of all such integers $`i`$ and, for each integer $`k1`$, denote by $`i_k`$ the $`k`$-th element of $`I`$. We claim that the sequences $`(𝐲_k)_{k1}`$ and $`(Y_k)_{k1}`$ given by

$$𝐲_k=𝐱_{i_k}\text{and}Y_k=X_{i_k},(k1),$$

have the required properties.

To prove this, fix any index $`k1`$ and put $`i=i_k`$. The inequality (5.2) gives

$$L_{i1}X_i^{1/\gamma }\text{and}L_iX_{i+1}^{1/\gamma }.$$

Then, by part (i) of Lemma 3.1 and the fact that $`det(𝐱_i)`$ is a non-zero integer, we find

$$1|det(𝐱_i)|X_iL_iX_iX_{i+1}^{1/\gamma }$$
$`(5.3)`$
and so $`X_{i+1}X_i^\gamma `$. Similarly, using part (ii) of Lemma 3.1 and the fact that the integral points $`𝐱_{i1}`$, $`𝐱_i`$ and $`𝐱_{i+1}`$ are linearly independent over $`𝐐`$, we then find

$$1|det(𝐱_{i1},𝐱_i,𝐱_{i+1})|X_{i+1}L_iL_{i1}X_{i+1}^{1/\gamma ^2}L_{i1}X_{i+1}^{1/\gamma ^2}X_i^{1/\gamma }1.$$

From these last estimates, we deduce

$$X_{i+1}X_i^\gamma ,L_{i1}X_i^{1/\gamma }\text{and}L_iX_{i+1}^{1/\gamma }X_i^1.$$
$`(5.4)`$
In particular, since $`𝐱_i=𝐲_k`$ and $`X_i=Y_k`$, we obtain, from (5.3) and (5.4),

$$𝐲_kY_k,L(𝐲_k)Y_k^1\text{and}1|det(𝐲_k)|1.$$
$`(5.5)`$
Now, let $`V=𝐱_i,𝐱_{i+1}_𝐐`$ and let $`ji+1`$ be the largest integer such that $`V`$ contains $`𝐱_i,𝐱_{i+1},\mathrm{},𝐱_j`$. Then we have

$$V=𝐱_i,𝐱_{i+1}_𝐐=𝐱_{j1},𝐱_j_𝐐$$
$`(5.6)`$
and, since $`𝐱_{j+1}V`$, we deduce that the points $`𝐱_{j1}`$, $`𝐱_j`$ and $`𝐱_{j+1}`$ are linearly independent over $`𝐐`$. Moreover, since $`𝐱_i,\mathrm{},𝐱_jV`$, the index $`j`$ is the smallest integer $`>i`$ with the latter property. Therefore, we have $`j=i_{k+1}`$. Consequently the estimates (5.4) also apply with $`i`$ replaced by $`j`$. Lemma 4.1 then gives

$$H\left(𝐱_i,𝐱_{i+1}_𝐐\right)X_{i+1}L_iX_i^{1/\gamma }\text{and}H\left(𝐱_{j1},𝐱_j_𝐐\right)X_jL_{j1}X_j^{1/\gamma ^2}.$$

By virtue of (5.6), this implies $`X_jX_i^\gamma `$ and, since $`X_i=Y_k`$ and $`X_j=Y_{k+1}`$, we conclude that

$$Y_{k+1}Y_k^\gamma .$$
$`(5.7)`$
The preceding argument also shows that $`𝐲_k,𝐲_{k+1}_𝐐`$ contains $`𝐱_{i+1}`$ and $`𝐱_{j1}`$. Therefore, $`𝐲_{k+1},𝐲_{k+2}_𝐐`$ contains $`𝐱_{j+1}`$ and so $`𝐲_k,𝐲_{k+1},𝐲_{k+2}_𝐐`$ contains the three linearly independent points $`𝐱_{j1}`$, $`𝐱_j`$ and $`𝐱_{j+1}`$. This implies that $`𝐲_k`$, $`𝐲_{k+1}`$ and $`𝐲_{k+2}`$ are linearly independent over $`𝐐`$ and, combining Lemma 3.1 (ii) with (5.5) and (5.7), we find

$$1|det(𝐲_k,𝐲_{k+1},𝐲_{k+2})|𝐲_{k+2}L(𝐲_k)L(𝐲_{k+1})Y_{k+2}Y_k^1Y_{k+1}^11.$$

This estimate together with (5.5) and (5.7) proves that the sequences $`(𝐲_k)_{k1}`$ and $`(Y_k)_{k1}`$ have all the required properties.

###### Corollary 5.2

Let $`\xi `$ be an extremal real number and let $`(𝐲_k)_{k1}`$ and $`(Y_k)_{k1}`$ be as in the statement of Theorem 5.1. Then, for any sufficiently large integer $`k3`$, the point $`𝐲_{k+1}`$ is a non-zero rational multiple of $`[𝐲_k,𝐲_k,𝐲_{k2}]`$.

Proof. Let $`k4`$ be an integer and let $`𝐰=[𝐲_k,𝐲_k,𝐲_{k+1}]`$. By Lemma 2.1 (i), we have

$$det(𝐰)=det(𝐲_k)^2det(𝐲_{k+1}).$$

Thus $`det(𝐰)0`$ and consequently, $`𝐰`$ is a non-zero point of $`𝐙^3`$. Using Lemma 3.1 (iii), we find

$$𝐰Y_k^2Y_{k+1}^1Y_{k2}\text{and}L(𝐰)Y_{k+1}Y_k^2Y_{k2}^1.$$

Then, applying Lemma 3.1 (ii), we obtain

$$\begin{array}{cc}\hfill |det(𝐰,𝐲_{k3},𝐲_{k2})|& Y_{k2}Y_{k2}^1Y_{k3}^1Y_{k3}^1,\hfill \\ \hfill |det(𝐰,𝐲_{k2},𝐲_{k1})|& Y_{k1}Y_{k2}^2Y_{k3}^{1/\gamma }.\hfill \end{array}$$

But the above two determinants are integers. So, if $`k`$ is sufficiently large, they must be zero. Since $`𝐲_{k3}`$, $`𝐲_{k2}`$ and $`𝐲_{k1}`$ are linearly independent over $`𝐐`$, this then forces $`𝐰`$ to be a rational multiple of $`𝐲_{k2}`$. Since, by Lemma 2.1 (iii), we have

$$[𝐲_k,𝐲_k,𝐰]=det(𝐲_k)^2𝐲_{k+1},$$

we conclude that $`𝐲_{k+1}`$ is a rational multiple of $`[𝐲_k,𝐲_k,𝐲_{k2}]`$.

As we will see in the next section, this corollary provides a way to construct extremal numbers. It also provides the following formula.

###### Corollary 5.3

In the notation of Corollary 5.2, we have

$$det(𝐲_{k2},𝐲_{k1},𝐲_k)𝐲_{k+1}=det(𝐲_{k2},𝐲_{k1},𝐲_{k+1})𝐲_k+det(𝐲_{k1},𝐲_k,𝐲_{k+1})𝐲_{k2}$$

for any sufficiently large value of $`k`$.

Proof. By Lemma 2.1 (ii), we have $`det(𝐲_{k2},𝐲_k,[𝐲_k,𝐲_k,𝐲_{k2}])=0`$ and this, by Corollary 5.2, implies that the points $`𝐲_{k2}`$, $`𝐲_k`$ and $`𝐲_{k+1}`$ are linearly dependent over $`𝐐`$. Writing $`𝐲_{k+1}=a𝐲_k+b𝐲_{k2}`$ with unknown rational coefficients $`a`$ and $`b`$ and using the multilinearity of the determinant, we get $`det(𝐲_{k2},𝐲_{k1},𝐲_{k+1})=adet(𝐲_{k2},𝐲_{k1},𝐲_k)`$ and $`det(𝐲_{k1},𝐲_k,𝐲_{k+1})=bdet(𝐲_{k1},𝐲_k,𝐲_{k2})`$ and the formula follows.

###### Corollary 5.4

The set $``$ of extremal real numbers is at most countable.

Proof. We prove this by constructing an injective map $`\phi :(𝐙^3)^3`$ as follows. For each extremal real number $`\xi `$, we choose a corresponding sequence $`(𝐲_k)_{k1}`$ as in Theorem 5.1. Then, in accordance with Corollary 5.3, we choose an index $`i3`$ such that, for each $`ki`$, the point $`𝐲_{k+1}`$ is a non-zero rational multiple of $`[𝐲_k,𝐲_k,𝐲_{k2}]`$, and we put $`\phi (\xi )=(𝐲_{i2},𝐲_{i1},𝐲_i)`$. Since the knowledge of $`𝐲_{i2}`$, $`𝐲_{i1}`$, $`𝐲_i`$ determines all points $`𝐲_k`$ with $`ki2`$ up to a non-zero rational multiple, it determines uniquely the corresponding ratios $`y_{k,1}/y_{k,0}`$ and also their limit $`\xi =lim_k\mathrm{}y_{k,1}/y_{k,0}`$. Thus, $`\phi `$ is injective.

In order to complete the proof of Theorem 1.1, it remains to show that the set of extremal numbers is infinite. This is the object of the next section.

6. Construction of extremal numbers

As in Section 2, we identify each point $`(y_0,y_1,y_2)`$ of $`𝐙^3`$ with the corresponding symmetric matrix with integral coefficients $`\left(\begin{array}{cc}y_0& y_1\\ y_1& y_2\end{array}\right)`$. Before stating the main result of this section, we first prove a lemma.

###### Lemma 6.1

Let $`A`$, $`B`$ be non-commuting symmetric matrices in $`\mathrm{GL}_2(𝐙)`$. Consider the sequence $`(𝐲_k)_{k1}`$ constructed recursively by putting

$$𝐲_1=B^1,𝐲_0=I,𝐲_1=A\text{and}𝐲_k=[𝐲_{k1},𝐲_{k1},𝐲_{k3}]\text{for }k2\text{,}$$
$`(6.1)`$
where $`I`$ denotes the identity $`2\times 2`$ matrix. Then, for each integer $`k1`$, we have

$$|det(𝐲_k)|=1\text{and}|det(𝐲_k,𝐲_{k+1},𝐲_{k+2})|=|\mathrm{trace}(JAB)|0$$
$`(6.2)`$
where $`J`$ is as defined in Section 2. Moreover, for $`k1`$, we have the alternative recurrence relation

$$𝐲_k=\pm 𝐲_{k1}S𝐲_{k2}\text{where}S=\{\begin{array}{cc}AB\hfill & \text{if }k\text{ is odd,}\hfill \\ BA\hfill & \text{if }k\text{ is even.}\hfill \end{array}$$
$`(6.3)`$
Proof. For $`k2`$, the recurrence relation in (6.1) implies, using Lemma 2.1 (i) and (ii),

$$\begin{array}{cc}\hfill det(𝐲_k)& =det(𝐲_{k1})^2det(𝐲_{k3}),\hfill \\ \hfill det(𝐲_{k2},𝐲_{k1},𝐲_k)& =det(𝐲_{k1})det(𝐲_{k3},𝐲_{k2},𝐲_{k1}).\hfill \end{array}$$

Since $`A`$ and $`B`$ belong to $`\mathrm{GL}_2(𝐙)`$, we also have $`det(𝐲_k)=\pm 1`$ for $`k=1,0,1`$, and therefore we deduce, by recurrence, that, for all $`k1`$, we have

$$|det(𝐲_k)|=1\text{and}|det(𝐲_k,𝐲_{k+1},𝐲_{k+2})|=|det(𝐲_1,𝐲_0,𝐲_1)|.$$

Moreover, using (2.1), we find $`det(𝐲_1,𝐲_0,𝐲_1)=det(B)\mathrm{trace}(JAB)0`$ since $`ABBA`$. This proves (6.2).

Similarly, a short computation shows that (6.3) holds for $`k=1`$ and, assuming that it holds for some integer $`k1`$, we find, using (2.2),

$$𝐲_{k+1}=𝐲_kJ𝐲_{k2}J𝐲_k=\pm 𝐲_{k1}S𝐲_{k2}J𝐲_{k2}J𝐲_k=\pm 𝐲_{k1}S𝐲_k$$

which, by taking transposes, implies $`𝐲_{k+1}=\pm 𝐲_k{}_{}{}^{t}S𝐲_{k1}`$. Thus (6.3) holds for all $`k1`$.

###### Theorem 6.2

Let $`A`$, $`B`$ and the sequence $`(𝐲_k)_{k1}`$ be as in Lemma 6.1. Assume that all coefficients of $`A`$ are non-negative and that all those of $`AB`$ are positive. Then, upon writing $`𝐲_k=(y_{k,0},y_{k,1},y_{k,2})`$, we have that $`y_{k,0}`$ is a non-zero integer for all $`k2`$ and that the limit

$$\xi =\underset{k\mathrm{}}{lim}y_{k,1}/y_{k,0}$$

exists and is an extremal real number. Moreover, letting $`Y_k=𝐲_k`$ for all $`k`$, the sequences $`(𝐲_k)_{k3}`$ and $`(Y_k)_{k3}`$ share, with respect to $`\xi `$, all the properties stated in Theorem 5.1.

Proof. By virtue of (6.2), the sequence $`(𝐲_k)_{k1}`$ satisfies the last two conditions of Theorem 5.1. Since $`𝐲_0`$ and $`𝐲_1`$ have non-negative coefficients, since their rows and columns are all non-zero, and since $`AB`$ and its transpose $`BA`$ have positive coefficients, the recurrence relation (6.3) implies that, for $`k2`$, the point $`𝐲_k`$ has non-zero coefficients, all of the same sign. Thus, the integer $`Y_k=𝐲_k`$ is positive for $`k2`$ and there exists a constant $`c_1>1`$ such that, for all $`k4`$, we have

$$Y_{k2}Y_{k1}<Y_kc_1Y_{k2}Y_{k1}.$$
$`(6.4)`$
In particular, the sequence $`(Y_k)_{k3}`$ is unbounded and monotone increasing. Similarly, we find that $`lim_k\mathrm{}y_{k,0}=\mathrm{}`$.

For each $`k2`$, define $`I_k`$ to be the interval of $`𝐑`$ with end points $`y_{k,1}/y_{k,0}`$ and $`y_{k,2}/y_{k,1}`$. For $`k3`$, the relation (6.3) shows that there are positive integers $`r,s,t,u`$ depending on $`k`$ such that

$$\left(\begin{array}{cc}y_{k,0}& y_{k,1}\\ y_{k,1}& y_{k,2}\end{array}\right)=\pm \left(\begin{array}{cc}y_{k1,0}& y_{k1,1}\\ y_{k1,1}& y_{k1,2}\end{array}\right)\left(\begin{array}{cc}r& s\\ t& u\end{array}\right)$$

and consequently we find

$$\frac{y_{k,1}}{y_{k,0}}=\frac{ry_{k1,1}+ty_{k1,2}}{ry_{k1,0}+ty_{k1,1}}I_{k1}\text{and}\frac{y_{k,2}}{y_{k,1}}=\frac{sy_{k1,1}+uy_{k1,2}}{sy_{k1,0}+uy_{k1,1}}I_{k1}.$$

Thus, the intervals $`I_k`$ with $`k2`$ form a non-increasing sequence $`I_2I_3I_4\mathrm{}`$ and, since the length of $`I_k`$ is $`1/(y_{k,0}y_{k,1})`$ tending to zero as $`k`$ tends to infinity, their intersection is reduced to one point $`\xi `$ with $`\xi >0`$. The ratios $`y_{k,1}/y_{k,0}`$ and $`y_{k,2}/y_{k,1}`$ therefore converge to that number $`\xi `$ and, since $`\xi 0`$, we get

$$|y_{k,0}||y_{k,1}||y_{k,2}|Y_k.$$

In particular, the length of $`I_k`$ satisfies $`|I_k|=(y_{k,0}y_{k,1})^1Y_k^2`$ and so we find

$$L(𝐲_k)\underset{j=0,1}{\mathrm{max}}|y_{k,j}\xi y_{k,j+1}|Y_k\underset{j=0,1}{\mathrm{max}}\left|\xi \frac{y_{k,j+1}}{y_{k,j}}\right|Y_k|I_k|Y_k^1.$$

To conclude from Theorem 5.1 that $`\xi `$ is extremal, it remains only to show that $`Y_kY_{k1}^\gamma `$. To this end, we observe that, for $`k4`$, the real number $`q_k=Y_kY_{k1}^\gamma `$ satisfies, by (6.4),

$$q_{k1}^{1/\gamma }q_kc_1q_{k1}^{1/\gamma }.$$

By recurrence on $`k`$, this implies that, for all $`k3`$, we have $`c_2^{1/\gamma }q_kc_2`$ for the constant $`c_2=\mathrm{max}\{c_1^\gamma ,q_3,q_3^\gamma \}`$ and so $`c_2^{1/\gamma }Y_{k1}^\gamma Y_kc_2Y_{k1}^\gamma `$ as required.

Let $`E=\{a,b\}`$ be an alphabet of two letters and let $`E^{}`$ denote the monoid of words on $`E`$ for the product given by concatenation. The Fibonacci sequence in $`E^{}`$ is the sequence of words $`(w_i)_{i0}`$ defined recursively by

$$w_0=b,w_1=a\text{and}w_i=w_{i1}w_{i2},(i2)$$

(see Example 1.3.6 of ). Since, for every $`i1`$, the word $`w_i`$ is a prefix of $`w_{i+1}`$, this sequence converges to an infinite word $`w=abaabab\mathrm{}`$ called the Fibonacci word on $`\{a,b\}`$. We now turn to Fibonacci continued fractions.

###### Corollary 6.3

Let $`a`$ and $`b`$ be distinct positive integers and let

$$\xi _{a,b}=[0,a,b,a,a,b,a,\mathrm{}]=1/(a+1/(b+\mathrm{}))$$

be the real number whose sequence of partial quotients starts with $`0`$ followed by the elements of the Fibonacci word on $`\{a,b\}`$. Then $`\xi _{a,b}`$ is an extremal real number. It is the special case of the construction of Theorem 6.2 corresponding to the choice of

$$A=\left(\begin{array}{cc}a& 1\\ 1& 0\end{array}\right)\text{and}B=\left(\begin{array}{cc}b& 1\\ 1& 0\end{array}\right).$$

Note that, if $`\xi `$ is an extremal real number, then $`1/\xi `$ is also extremal. So, the corollary implies that $`1/\xi _{a,b}=[a,b,a,a,b,\mathrm{}]`$ as well is extremal. More generally, one finds that, if $`\xi `$ is extremal, then the ratio $`(r\xi +s)/(t\xi +u)`$ is also extremal for any choice of rational numbers $`r,s,t,u`$ with $`rust0`$.

Proof. Recall that, if a real number $`\xi `$ has an infinite continued fraction expansion $`\xi =[a_0,a_1,a_2,\mathrm{}]`$ and if, for each $`i0`$, we write its convergent $`[a_0,a_1,\mathrm{},a_i]`$ in the form $`p_i/q_i`$ with relatively prime integers $`p_i`$ and $`q_i>0`$, then the sequences $`(p_i)_{i0}`$ and $`(q_i)_{i0}`$ satisfy recurrence relations which, in matrix form, can be written for $`i2`$ as

$$\left(\begin{array}{cc}q_i& q_{i1}\\ p_i& p_{i1}\end{array}\right)=\left(\begin{array}{cc}q_{i1}& q_{i2}\\ p_{i1}& p_{i2}\end{array}\right)\left(\begin{array}{cc}a_i& 1\\ 1& 0\end{array}\right)$$

(see for example Chapter I of ). When $`a_0=0`$, this gives

$$\left(\begin{array}{cc}q_i& q_{i1}\\ p_i& p_{i1}\end{array}\right)=\left(\begin{array}{cc}a_1& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}a_2& 1\\ 1& 0\end{array}\right)\mathrm{}\left(\begin{array}{cc}a_i& 1\\ 1& 0\end{array}\right).$$

Now, consider the morphism of monoids $`\phi :E^{}\mathrm{GL}_2(𝐙)`$ determined by the conditions

$$\phi (a)=A\text{and}\phi (b)=B.$$

For each integer $`k0`$, let $`m_k`$ denote the word $`w_{k+2}`$ minus its last two letters and let $`𝐱_k=\phi (m_k)`$. Then, we have $`m_0=1`$ (the empty word), $`m_1=a`$ and, for $`k2`$,

$$m_k=m_{k1}sm_{k2}\text{where}s=\{\begin{array}{cc}ab\hfill & \text{if }k\text{ is odd,}\hfill \\ ba\hfill & \text{if }k\text{ is even.}\hfill \end{array}$$

In terms of the matrices $`𝐱_k`$, this implies $`𝐱_0=I`$, $`𝐱_1=A`$ and, for $`k2`$,

$$𝐱_k=𝐱_{k1}S𝐱_{k2}\text{where}S=\{\begin{array}{cc}AB\hfill & \text{if }k\text{ is odd,}\hfill \\ BA\hfill & \text{if }k\text{ is even.}\hfill \end{array}$$

Therefore, if $`(𝐲_k)_{k1}`$ denotes the sequence furnished by Lemma 6.1 for the above choice of $`A`$ and $`B`$, we deduce from (6.3) that $`𝐲_k=\pm 𝐱_k`$ for any $`k0`$. In particular, $`𝐱_k`$ is a symmetric matrix and, with the usual notation for coordinates, we have

$$\frac{y_{k,1}}{y_{k,0}}=\frac{x_{k,1}}{x_{k,0}}$$

for all $`k1`$. On one hand, Theorem 5.1 tells us that the above ratio converges to an extremal real number $`\xi `$ as $`k`$ tends to infinity. On the other hand, since $`m_k`$ is a prefix of the Fibonacci word $`w`$ and since $`𝐱_k=\phi (m_k)`$, the observation made at the beginning of the proof shows that the same ratio is the convergent of $`\xi _{a,b}`$ whose partial quotients are $`0`$ followed by the elements of $`m_k`$. Therefore $`\xi _{a,b}=\xi `$ is extremal.

7. Approximation by rational numbers

In this section, we prove Theorem 1.3, as a consequence of the following result.

###### Proposition 7.1

Let $`\xi `$ be an extremal real number. There exists a positive constant $`s`$ such that, for any sufficiently large real number $`X`$, the first minimum of the convex body

$$𝒞(X)=\{(x_0,x_1)𝐑^2;|x_0|X\text{and}|x_0\xi x_1|X^1\}$$

is bounded below by $`(\mathrm{log}X)^s`$.

Proof. Let $`(Y_k)_{k1}`$ and $`(𝐲_k)_{k1}`$ be the sequences given by Theorem 5.1. Let $`X`$ be a real number with $`XY_2`$, let $`\lambda =\lambda (X)`$ be the first minimum of $`𝒞(X)`$, and let $`(x_0,x_1)`$ be a point of $`𝐙^2`$ which realizes this minimum, so that we have $`|x_0|\lambda X`$ and $`|x_0\xi x_1|\lambda X^1`$. We choose an index $`k2`$ such that $`Y_kXY_{k+1}`$ and consider the point $`(z_0,z_1)`$ of $`𝐙^2`$ given by

$$(z_0,z_1)=(x_0,x_1)J𝐲_{k+1}J𝐲_{k1}$$

in the notation of Section 2. Writing $`𝐲_{k1}=(y_0^{},y_1^{},y_2^{})`$ and $`𝐲_{k+1}=(y_0^{},y_1^{},y_2^{})`$, and using the estimates of Theorem 5.1 together with Lemma 3.1 (i), we find

$$\begin{array}{cc}\hfill |z_0|& =\begin{array}{cc}x_0& x_1\\ \\ \left|\begin{array}{cc}y_0^{}& y_1^{}\\ y_0^{}& y_1^{}\end{array}\right|& \left|\begin{array}{cc}y_0^{}& y_1^{}\\ y_1^{}& y_2^{}\end{array}\right|\end{array}=\begin{array}{cc}x_0& x_0\xi x_1\\ \\ \left|\begin{array}{cc}y_0^{}& y_1^{}\\ y_0^{}& y_1^{}\end{array}\right|& \left|\begin{array}{cc}y_0^{}& y_1^{}\\ y_0^{}\xi y_1^{}& y_1^{}\xi y_2^{}\end{array}\right|\end{array}\hfill \\ & \\ & \frac{|x_0|}{Y_k}+Y_k|x_0\xi x_1|\frac{\lambda X}{Y_k}\hfill \end{array}$$

and

$$\begin{array}{cc}\hfill |z_0\xi z_1|& =\begin{array}{cc}x_0& x_1\\ \\ \left|\begin{array}{cc}y_0^{}\xi y_1^{}& y_1^{}\xi y_2^{}\\ y_0^{}& y_1^{}\end{array}\right|& \left|\begin{array}{cc}y_0^{}\xi y_1^{}& y_1^{}\xi y_2^{}\\ y_1^{}& y_2^{}\end{array}\right|\end{array}\hfill \\ & \\ & =\begin{array}{cc}x_0& x_0\xi x_1\\ \\ \left|\begin{array}{cc}y_0^{}\xi y_1^{}& y_1^{}\xi y_2^{}\\ y_0^{}& y_1^{}\end{array}\right|& \left|\begin{array}{cc}y_0^{}\xi y_1^{}& y_1^{}\xi y_2^{}\\ y_0^{}\xi y_1^{}& y_1^{}\xi y_2^{}\end{array}\right|\end{array}\hfill \\ & \\ & \frac{|x_0|}{Y_{k1}Y_{k+1}}+Y_k|x_0\xi x_1|\frac{\lambda Y_k}{X}.\hfill \end{array}$$

Thus, if we put $`Z=2X/Y_k`$, these relations imply that the first minimum $`\lambda (Z)`$ of $`𝒞(Z)`$ satisfies $`\lambda (Z)c\lambda (X)`$ for some constant $`c>0`$ which is independent of $`k`$ and $`X`$. We also note that $`ZX^{1/\gamma ^2}`$ and thus $`ZX^{1/2}`$ if $`X`$ is sufficiently large, say $`XX_0Y_2`$. Then, by choosing $`s>0`$ with $`2^{s1}c`$, we get, for $`XX_0`$,

$$\lambda (X)(\mathrm{log}X)^s2\lambda (Z)(\mathrm{log}Z)^s$$

where $`Z`$ is a real number in the interval $`2ZX^{1/2}`$. Since $`\lambda (X)(\mathrm{log}X)^s`$ is bounded below by some positive constant for $`2XX_0^{1/2}`$, this shows that $`\lambda (X)(\mathrm{log}X)^s`$ tends to infinity with $`X`$.

Proof of Theorem 1.3. Let $`s`$ be as in the proposition and assume that a rational number $`\alpha =p/q`$ with denominator $`q3`$ satisfies $`|\xi p/q|q^2(\mathrm{log}q)^{2s2}`$. Then, by putting $`X=q(\mathrm{log}q)^{s+1}`$ and $`\lambda =(\mathrm{log}q)^{s1}`$, we get $`|q|\lambda X`$ and $`|q\xi p|\lambda /X`$ so that the first minimum of $`𝒞(X)`$ is $`\lambda `$. By the proposition, this implies $`\lambda (\mathrm{log}X)^s`$ if $`q`$ is sufficiently large. This in turn forces an upper bound on $`q`$ and therefore on $`H(p/q)`$. Thus Theorem 1.3 holds with $`t=2s+2`$.

8. Approximation by quadratic real numbers

In this section, we prove Theorems 1.2 and 1.4 as consequences of the following result.

###### Proposition 8.1

Let $`\xi `$ be an extremal real number and let the sequences $`(Y_k)_{k1}`$ and $`(𝐲_k)_{k1}`$ be as in Theorem 5.1. For each $`k1`$, we define

$$Q_k(T)=\left|\begin{array}{ccc}1& T& T^2\\ y_{k,0}& y_{k,1}& y_{k,2}\\ y_{k+1,0}& y_{k+1,1}& y_{k+1,2}\end{array}\right|.$$

Then, $`Q_k`$ is a polynomial with integral coefficients which, for all sufficiently large values of $`k`$, has degree $`2`$ and satisfies

$$H(Q_k)|Q_k^{}(\xi )|Y_{k1}\text{and}|Q_k(\xi )|Y_{k+2}^1.$$

Proof. By Lemma 3.1 (i), we have

$$H(Q_k)=\begin{array}{ccc}y_{k,0}& y_{k,1}& y_{k,2}\\ y_{k+1,0}& y_{k+1,1}& y_{k+1,2}\end{array}Y_{k+1}L(𝐲_k)+Y_kL(𝐲_{k+1})Y_{k1}.$$

All the remaining estimates assume that $`k`$ is sufficiently large. Since

$$Q_k^{\prime \prime }=2\left|\begin{array}{cc}y_{k,0}& y_{k,1}y_{k,0}\xi \\ y_{k+1,0}& y_{k+1,1}y_{k+1,0}\xi \end{array}\right|=2y_{k+1,0}(y_{k,0}\xi y_{k,1})+𝒪(Y_{k1}^1),$$

we get $`Q_k^{\prime \prime }0`$ for $`k1`$ using the lower bound on $`|y_{k,0}\xi y_{k,1}|`$ provided by Theorem 1.3. Similarly, we have

$$Q_k^{}(\xi )=\left|\begin{array}{cc}y_{k,0}& y_{k,2}2y_{k,1}\xi +y_{k,0}\xi ^2\\ y_{k+1,0}& y_{k+1,2}2y_{k+1,1}\xi +y_{k+1,0}\xi ^2\end{array}\right|.$$
$`(8.1)`$
Note that, for any $`j1`$,

$$y_{j,0}(y_{j,2}2y_{j,1}\xi +y_{j,0}\xi ^2)=det(𝐲_j)+(y_{j,1}y_{j,0}\xi )^2=det(𝐲_j)+𝒪(Y_j^2),$$

and therefore

$$|y_{j,2}2y_{j,1}\xi +y_{j,0}\xi ^2|Y_j^1.$$

Using this to estimate (8.1), we find

$$|Q_k^{}(\xi )|=|y_{k+1,0}||y_{k,2}2y_{k,1}\xi +y_{k,0}\xi ^2|+𝒪(Y_{k1}^1)Y_{k1}$$

and, since $`H(Q_k)Y_{k1}`$, we get $`H(Q_k)Y_{k1}`$. Finally, write

$$y_{k+2,0}(1,\xi ,\xi ^2)=𝐲_{k+2}+𝐳$$

where $`𝐳𝐑^3`$ has norm $`𝐳Y_{k+2}^1`$. By multilinearity of the determinant, this gives

$$y_{k+2,0}Q_k(\xi )=det(𝐲_{k+2},𝐲_k,𝐲_{k+1})+det(𝐳,𝐲_k,𝐲_{k+1}).$$

Since

$$|det(𝐳,𝐲_k,𝐲_{k+1})|3𝐳H(Q_k)Y_{k1}/Y_{k+2},$$

we conclude that

$$|Q_k(\xi )|=|y_{k+2,0}|^1\left(|det(𝐲_k,𝐲_{k+1},𝐲_{k+2})|+𝒪(Y_{k1}/Y_{k+2})\right)Y_{k+2}^1.\text{ }\text{  }\text{ }$$

Proof of Theorem 1.2. For any real number $`X`$ with $`XH(Q_1)`$, we choose an index $`k1`$ such that $`H(Q_k)XH(Q_{k+1})`$. Then upon writing $`Q_k(T)=x_0T^2+x_1T+x_2`$, we find that $`(x_0,x_1,x_2)`$ is a non-zero integral point with $`|x_0|X`$, $`|x_1|X`$ and

$$|x_0\xi ^2+x_1\xi +x_2|=|Q_k(\xi )|Y_{k+2}^1H(Q_{k+1})^{\gamma ^2}X^{\gamma ^2}.\text{ }\text{  }\text{ }$$

As for Theorem 1.4, it follows immediately from the following more precise result.

###### Theorem 8.2

Let $`\xi `$ be an extremal real number. Then there exist an integer $`k_01`$ and positive constants $`c_1`$, $`c_2`$ such that, for any $`kk_0`$, the polynomial $`Q_k`$ of Lemma 8.1 is irreducible over $`𝐐`$ of degree $`2`$ and admits exactly one root $`\alpha _k`$ with

$$c_1H(\alpha _k)^{2\gamma ^2}|\xi \alpha _k|c_2H(\alpha _k)^{2\gamma ^2}.$$
$`(8.2)`$
There also exists a constant $`c_3>0`$ such that, for any algebraic number $`\alpha 𝐂`$ of degree $`2`$ over $`𝐐`$, distinct from all $`\alpha _k`$ with $`kk_0`$, we have

$$|\xi \alpha |c_3H(\alpha )^4.$$

Proof. By Proposition 8.1, the polynomial $`Q_k`$ has degree $`2`$ for $`k1`$. For these values of $`k`$, we may factor it as $`Q_k(T)=a_k(T\alpha _k)(T\beta _k)`$ with the roots $`\alpha _k`$ and $`\beta _k`$ ordered so that $`|\xi \alpha _k||\xi \beta _k|`$. Then the logarithmic derivative of $`Q_k`$ satisfies

$$\frac{1}{|\xi \alpha _k|}\frac{1}{|\xi \beta _k|}\frac{|Q_k^{}(\xi )|}{|Q_k(\xi )|}=\left|\frac{1}{\xi \alpha _k}+\frac{1}{\xi \beta _k}\right|\frac{2}{|\xi \alpha _k|}.$$
$`(8.3)`$
Using the estimates of the proposition, this gives

$$|\xi \alpha _k|2\frac{|Q_k(\xi )|}{|Q_k^{}(\xi )|}H(Q_k)^{2\gamma ^2}H(\alpha _k)^{2\gamma ^2}.$$
$`(8.4)`$
Since $`H(Q_k)`$ tends to infinity with $`k`$, we deduce that $`\alpha _k`$ tends to $`\xi `$ as $`k`$ tends to infinity and, by comparison with Theorem 1.3, that $`\alpha _k`$ is irrational for all sufficiently large values of $`k`$. Thus, for $`k1`$, the polynomial $`Q_k`$ is irreducible over $`𝐐`$. Its discriminant $`a_k^2(\alpha _k\beta _k)^2`$ is then a non-zero integer, and so we find

$$|\alpha _k\beta _k||a_k|^1H(Q_k)^1.$$

In particular, we have $`|\alpha _k\beta _k|3|\xi \alpha _k|`$ for any sufficiently large value of $`k`$, and (8.3) becomes

$$\frac{1}{2|\xi \alpha _k|}\frac{|Q_k^{}(\xi )|}{|Q_k(\xi )|}\frac{2}{|\xi \alpha _k|}.$$

On the other hand, the greatest common divisor of the coefficients of $`Q_k`$ divides the non-zero integer $`det(𝐲_{k1},𝐲_k,𝐲_{k+1})`$ and, since the absolute value of this integer is bounded above by a constant that is independent of $`k`$, we also have $`H(\alpha _k)H(Q_k)`$ for $`k1`$. Then, going back to (8.4), we get refined estimates of the form (8.2) valid say for all $`kk_0`$, where $`c_1`$, $`c_2`$ are positive constants and $`k_02`$ is a fixed integer. This proves the first part of the Theorem.

Now, let $`\alpha 𝐂`$ be an algebraic number of degree $`2`$ with $`\alpha \alpha _k`$ for any index $`kk_0`$. For these values of $`k`$, Liouville’s inequality gives

$$|\alpha \alpha _k|c_3H(\alpha )^2H(\alpha _k)^2$$

where $`c_3>0`$ is an absolute constant (see for example Lemma 3 of ). Choose $`k`$ to be the smallest integer $`k_0`$ satisfying

$$H(\alpha )\left(\frac{c_3}{2c_2}\right)^{1/2}H(\alpha _k)^\gamma .$$

This choice of $`k`$ ensures that

$$|\xi \alpha _k|c_2H(\alpha _k)^{2\gamma ^2}\frac{c_3}{2}H(\alpha )^2H(\alpha _k)^2\frac{1}{2}|\alpha \alpha _k|,$$

and therefore that

$$|\xi \alpha |\frac{1}{2}|\alpha \alpha _k|H(\alpha )^2H(\alpha _k)^2.$$

If $`k>k_0`$, we also have $`H(\alpha )H(\alpha _{k1})^\gamma H(\alpha _k)`$, while if $`k=k_0`$, we have $`H(\alpha )1H(\alpha _k)`$. So, the above inequality leads to $`|\xi \alpha |H(\alpha )^4`$.

9. Approximation by cubic algebraic integers

For a real number $`x`$, let $`\left\{x\right\}=\mathrm{min}\{|xp|;p𝐙\}`$ denote the distance from $`x`$ to a nearest integer. In this section, we fix an extremal real number $`\xi `$ and corresponding sequences $`(Y_k)_{k1}`$ and $`(𝐲_k)_{k1}`$ satisfying the conditions of Theorem 5.1. In particular, we have $`\left\{y_{k,0}\xi ^j\right\}Y_k^1`$ for $`k1`$ and $`j=0,1,2`$. Here, we prove a lower bound for $`\left\{y_{k,0}\xi ^3\right\}`$ which implies the measure of approximation of Theorem 1.5 through the following proposition.

###### Proposition 9.1

Assume that there exist real numbers $`\delta `$ and $`c_1`$ with $`0\delta <1`$ and $`c_1>0`$ such that

$$\left\{y_{k,0}\xi ^3\right\}c_1Y_k^\delta $$
$`(9.1)`$
for all $`k1`$. Then there exists a constant $`c_2>0`$ such that, for any algebraic integer $`\alpha 𝐂`$ of degree $`3`$, we have

$$|\xi \alpha |c_2H(\alpha )^\theta $$

where $`\theta =(\gamma ^2+\delta /\gamma )/(1\delta )`$.

Proof. Let $`\alpha `$ be an algebraic integer of degree $`3`$ and let $`P(T)=T^3+pT^2+qT+r`$ be the product of its irreducible polynomial over $`𝐙`$ by the appropriate power of $`T`$ which makes it of degree $`3`$. Since $`H(P)=H(\alpha )`$ we find, for all $`k1`$,

$$\begin{array}{cc}\hfill \left\{y_{k,0}\xi ^3\right\}& |y_{k,0}P(\xi )|+|p|\left\{y_{k,0}\xi ^2\right\}+|q|\left\{y_{k,0}\xi \right\}\hfill \\ & c_3\left(Y_kH(\alpha )|\xi \alpha |+Y_k^1H(\alpha )\right)\hfill \end{array}$$

with a constant $`c_3>0`$ depending only on $`\xi `$. Choosing $`k`$ to be the smallest positive integer for which

$$H(\alpha )\frac{c_1}{2c_3}Y_k^{1\delta },$$

and using (9.1) this implies

$$|\xi \alpha |\frac{c_1}{2c_3}Y_k^{1\delta }H(\alpha )^1.$$

The conclusion follows since the above choice of $`k`$ implies $`Y_kH(\alpha )^{\gamma /(1\delta )}`$.

Numerical experiments done on Fibonacci continued fractions suggest that the values of $`\left\{y_{k,0}\xi ^3\right\}`$ are more or less uniformly distributed in the interval $`(0,1/2)`$. If true, this would imply a lower bound of the form $`\left\{y_{k,0}\xi ^3\right\}1/k(\mathrm{log}\mathrm{log}Y_k)^1`$. To be safer, we may hope for a lower bound of the type $`\left\{y_{k,0}\xi ^3\right\}Y_k^\delta `$ for any $`\delta >0`$. Then, for any $`ϵ>0`$, the above proposition gives a measure of approximation by algebraic integers $`\alpha `$ of degree $`3`$ of the form $`|\xi \alpha |cH(\alpha )^{\gamma ^2ϵ}`$ with a constant $`c`$ depending on $`\xi `$ and $`ϵ`$. Here we content ourselves with the following much weaker lower bound which, in the same way, implies Theorem 1.5.

###### Proposition 9.2

In the above notation, we have $`\left\{y_{k,0}\xi ^3\right\}Y_k^{1/\gamma ^3}`$ for all sufficiently large values of $`k`$ (to avoid those indices with $`y_{k,0}=0`$).

Proof. Fix an integer $`k4`$ and denote by $`y_{k,3}`$ the nearest integer to $`y_{k,2}\xi `$. Define

$$\stackrel{~}{𝐲}_k=(y_{k,1},y_{k,2},y_{k,3}),d_{k2}=det(𝐲_{k2},𝐲_{k1},𝐲_k)\text{and}\stackrel{~}{d}_{k2}=det(𝐲_{k2},𝐲_{k1},\stackrel{~}{𝐲}_k,).$$

According to Theorem 5.1, $`d_{k2}`$ belongs to a finite set of non-zero integers. Since $`\xi `$ is irrational and since $`\stackrel{~}{d}_{k2}`$ is an integer, we deduce that

$$1|d_{k2}\xi \stackrel{~}{d}_{k2}|.$$

On the other hand, using Proposition 8.1, we find

$$\begin{array}{cc}\hfill |d_{k2}\xi \stackrel{~}{d}_{k2}|& =|det(𝐲_{k2},𝐲_{k1},\xi 𝐲_k\stackrel{~}{𝐲}_k)|\hfill \\ & \xi 𝐲_k\stackrel{~}{𝐲}_kH(Q_{k2})\hfill \\ & \left(\left\{y_{k,2}\xi \right\}+Y_k^1\right)Y_{k3}.\hfill \end{array}$$

Combining these two estimates, we obtain $`\left\{y_{k,2}\xi \right\}Y_{k3}^1Y_k^{1/\gamma ^3}`$ for all sufficiently large $`k`$. The conclusion follows since $`|y_{k,0}\xi ^3y_{k,2}\xi |Y_k^1`$.

10. A partial converse

We conclude this paper with the following partial converse to Proposition 9.1, where the notation is as in §9.

###### Proposition 10.1

Let $`\mathrm{}`$ be the least common multiple of all integers of the form $`d_k=det(𝐲_k,𝐲_{k+1},𝐲_{k+2})`$ with $`k1`$, and assume that there exist infinitely many indices $`k1`$ satisfying

$$\left\{\frac{y_{k,0}\xi ^3}{\mathrm{}}\right\}c_1Y_k^\delta $$
$`(10.1)`$
for some real numbers $`c_1>0`$ and $`\delta >0`$. Then there exist a constant $`c_2>0`$ and infinitely many algebraic integers $`\alpha 𝐂`$ of degree $`3`$ for which

$$|\xi \alpha |c_2H(\alpha )^\theta $$

where $`\theta =(\gamma ^2\delta )/(1\delta )`$.

Our proof below goes by making explicit the arguments of Davenport and Schmidt in §2 of . To make the connection clearer for the interested reader, denote by $`𝒞_k`$ the convex body of $`𝐑^3`$ defined by

$$|x_0|Y_k^\gamma ,|x_0\xi x_1|Y_k^1\text{and}|x_0\xi ^2x_2|Y_k^1.$$

Then there exists a constant $`c>0`$ which is independent of $`k`$ such that $`𝐲_{k+1}c𝒞_k`$, $`𝐲_kc𝒞_k`$ and $`𝐲_{k1}cY_k^{1/\gamma ^2}𝒞_k`$ for $`k2`$. Since the volume of $`𝒞_k`$ is $`8Y_k^{1/\gamma ^2}`$ and since $`𝐲_{k+1}`$, $`𝐲_k`$ and $`𝐲_{k1}`$ are linearly independent over $`𝐐`$, this implies that the successive minima of $`𝒞_k`$ behave like $`1`$, $`1`$ and $`Y_k^{1/\gamma ^2}`$ and that they are essentially realized by these three points. Then, the successive minima of the polar convex body $`𝒞_k^{}`$ defined by

$$|x_0+x_1\xi +x_2\xi ^2|Y_k^\gamma ,|x_1|Y_k\text{and}|x_2|Y_k$$

behave like $`Y_k^{1/\gamma ^2}`$, $`1`$ and $`1`$ and, by identifying a point $`(x_0,x_1,x_2)𝐙^3`$ with the polynomial $`x_0+x_1T+x_2T^2𝐙[T]`$, these minima are essentially realized by the polynomials denoted $`B`$, $`A`$ and $`C`$ in the proof below.

Proof. Fix an index $`k2`$ for which (10.1) is satisfied. Define $`𝐓=(1,T,T^2)`$ and consider the three polynomials

$$A=det(𝐓,𝐲_{k1},𝐲_{k+1}),B=det(𝐓,𝐲_k,𝐲_{k+1})\text{and}C=det(𝐓,𝐲_{k1},𝐲_k).$$

We find

$$d_{k1}=y_{k,0}A(T)+y_{k1,0}B(T)+y_{k+1,0}C(T).$$

Thus, if we put $`m=\mathrm{}/d_{k1}`$, then, for a suitable choice of signs, the polynomial

$$P(T)=T^3\xi ^3\pm m\left\{\frac{y_{k,0}\xi ^3}{\mathrm{}}\right\}A(T)\pm m\left\{\frac{y_{k1,0}\xi ^3}{\mathrm{}}\right\}B(T)\pm m\left\{\frac{y_{k+1,0}\xi ^3}{\mathrm{}}\right\}C(T)$$

has integral coefficients and is monic of degree three. We claim that it satisfies also

$$|P^{}(\xi )|H(P)\text{and}|P(\xi )|H(P)^{\gamma /(1\delta )}$$

and that $`H(P)`$ tends to infinity with $`k`$. If we take this for granted, then the root $`\alpha `$ of $`P`$ which is closest to $`\xi `$ is an algebraic integer of degree $`3`$ with

$$|\xi \alpha |\frac{|P(\xi )|}{|P^{}(\xi )|}H(P)^\theta H(\alpha )^\theta $$

and the proposition is proved.

To establish the claim, we first note that, by (10.1) and Proposition 9.2, we have

$$Y_k^\delta \left\{\frac{y_{k,0}\xi ^3}{\mathrm{}}\right\}\frac{1}{\mathrm{}}\left\{y_{k,0}\xi ^3\right\}Y_k^{1/\gamma ^3},$$
$`(10.2)`$
and so $`\delta \gamma ^3`$ since $`k`$ may be arbitrarily large. We also observe that, in the notation of Proposition 8.1, we have $`B=Q_k`$ and $`C=Q_{k1}`$. Thus, with the appropriate choice of sign, we get

$$H\left(P\pm m\left\{\frac{y_{k,0}\xi ^3}{\mathrm{}}\right\}A\right)1+H(B)+H(C)Y_{k1}$$
$`(10.3)`$
and

$$|P(\xi )||A(\xi )|+|B(\xi )|+|C(\xi )||A(\xi )|+Y_{k+1}^1.$$
$`(10.4)`$
Moreover, by Corollary 5.3, the point $`d_{k1}𝐲_{k+2}d_k𝐲_{k1}`$ is an integral multiple of $`𝐲_{k+1}`$ and thus we find

$$d_{k1}Q_{k+1}+d_kA=det(𝐓,𝐲_{k+1},d_{k1}𝐲_{k+2}d_k𝐲_{k1})=0.$$

Since $`d_{k1}`$ and $`d_k`$ are non-zero integers of bounded absolute value, this implies, by Proposition 8.1, that

$$H(A)|A^{}(\xi )|Y_k\text{and}|A(\xi )|Y_{k+3}^1.$$
$`(10.5)`$
Combining (10.2), (10.3) and (10.5) we deduce

$$H(P)\left\{\frac{y_{k,0}\xi ^3}{\mathrm{}}\right\}Y_k|P^{}(\xi )|.$$

In particular this gives $`H(P)Y_k^{1\delta }`$ and so, using (10.4) and (10.5), we get

$$|P(\xi )|Y_k^\gamma H(P)^{\gamma /(1\delta )}$$

which completes the proof of the claim.

###### Corollary 10.2

Let the hypotheses be as in Theorem 6.2 and assume moreover that $`\mathrm{trace}(JAB)=\pm 1`$. Then the following conditions are equivalent:

(i) There exists $`\theta >\gamma ^2`$ such that the inequality $`|\xi \alpha |H(\alpha )^\theta `$ has infinitely many solutions in algebraic integers $`\alpha 𝐂`$ of degree $`3`$ over $`𝐐`$.

(ii) There exists $`\delta >0`$ such that the inequality $`\{y_{k,0}\xi ^3\}Y_k^\delta `$ holds for infinitely many indices $`k`$.

Proof. The fact that (i) implies (ii) follows from Proposition 9.1. The reverse implication follows from Proposition 10.1 granted that, for the sequence $`(𝐲_k)_{k1}`$ constructed by Lemma 6.1, we have $`|det(𝐲_k,𝐲_{k+1},𝐲_{k+2})|=1`$ for all $`k1`$.

In particular, the corollary applies to any Fibonacci continued fraction $`\xi _{a,b}`$ for which $`|ab|=1`$.

Acknowledgments. The author thanks Drew Vandeth for the reference  and for recognizing the continued fraction expansion of Corollary 6.3 as a Fibonacci word.

References

 J.-P. Allouche, J. L. Davison, M. Queffélec, L. Q. Zamboni, Transcendence of Sturmian or morphic continued fractions, J. Number Theory 91 (2001), 39–66.

 Y. Bugeaud, O. Teulié, Approximation d’un nombre réel par des nombres algébriques de degré donné, Acta Arith. 93 (2000), 77–86.

 J. W. S. Cassels, An introduction to Diophantine approximation, Cambridge Univ. Press, 1957.

 H. Davenport, W. M. Schmidt, Approximation to real numbers by quadratic irrationals, Acta Arith. 13 (1967), 169-176.

 H. Davenport, W. M. Schmidt, Approximation to real numbers by algebraic integers, Acta Arith. 15 (1969), 393–416.

 M. Laurent, D. Roy, Criteria of algebraic independence with multiplicities and interpolation determinants, Trans. Amer. Math. Soc. 351 (1999), 1845–1870.

 M. Lothaire, Combinatorics on words, Encyclopedia of mathematics and its applications, vol. 17, Addison-Wesley Pub. Co., 1983.

 D. Roy, Approximation simultanée d’un nombre et de son carré, C. R. Acad. Sc. Paris, Série I, 336 (2003), 1–6, arXiv:math.NT/0210395.

 W. M. Schmidt, Diophantine approximation, Lecture Note in Math., vol. 785, Springer-Verlag, 1980.

 E. Wirsing, Approximation mit algebraischen Zahlen beschränkten Grades, J. reine angew. Math. 206 (1961), 67-77.

Damien ROY Département de Mathématiques et de Statistiques Université d’Ottawa 585 King Edward Ottawa, Ontario K1N 6N5, Canada E-mail: droy@uottawa.ca http$`:`$//aix1.uottawa.ca/$``$droy/


text:
# Dynamical tunneling in molecules: role of the classical resonances and chaos

## Abstract

In this letter we study dynamical tunneling in highly excited symmetric molecules. The role of classical phase space structures like resonances and chaos on the tunneling splittings are illustrated using the water molecule as an example. It is argued that the enhancements in the splittings due to resonances (near-integrable phase space) and due to chaos (mixed phase space) are best understood away from the fluctuations associated with avoided crossings. In particular we provide an essential difference between the two mechanisms in terms of high order perturbation theory. The analysis, apart from testing the validity of a perturbative approach, suggests such systems as prime candidates for studying dynamical tunneling.

The concept of dynamical tunneling emerged more than two decades ago in the context of studying near degenerate vibrational states of symmetric molecules. As an example in the water molecule the symmetry of the OH bonds implies that a state $`(n_1=3,n_2=0)`$ which has three quanta of excitation in one of the OH bond ($`n_1=3`$) and zero quanta of excitation in the other equivalent OH bond ($`n_2=0`$) is degenerate with the state $`(0,3)`$. The two states are not coupled through any classical process i.e., a classical trajectory started with initial conditions corresponding to the state $`(3,0)`$ will remain localized indefinitely without ever evolving to the part of the phase space corresponding to the state $`(0,3)`$. However, Lawton and Childlawch  realized that a generalized form of tunneling mixes the two degenerate states giving rise to a characteristic splitting. In an influential work Davis and Hellerdavhel  refined and generalized the earlier observations and introduced the term dynamical tunneling to distinguish it from the usual tunneling through potential barriers in coordinate space. One of the important outcomes of the study was the suggestion that the underlying phase space of the system is the proper setting to understand dynamical tunneling.

In an apparently unrelated developement researchers studying the problem of intramolecular energy redistribution (IVR) discovered purely quantum energy flow between modes which would be otherwise uncoupled. A detailed analysis by Hutchinson, Sibert and Hyneshusihy  established that the mechanism for such clasically forbidden energy flow between degenerate vibrational modes arose from indirect state-to-state explorations involving a sequence of intermediate states. Important insights were provided by Stuchebrukhov and Marcusstumar  who established that one could view dynamical tunneling as a high order perturbative process involving a sequence of off-resonance virtual states (“vibrational superexchange”). Interestingly it was also shown that the perturbative approach to the tunneling splitting was related in a simple way to the usual semiclassical JWKB solution.

Despite this important prior work the explicit demonstration of the role of various classical phase space structures in dynamical tunneling has only recently been established. In particular the work of Bohigas, Tomsovic and Ullmocat  elucidating the role of chaos on dynamical tunneling has led to a resurgence of interest in the field. A key observation that emerged from numerous studiescats  of chaos-assisted tunneling (CAT) is that such a process necessarily requires atleast a three level mechanism. One of the hallmarks of CAT has to do with erratic fluctuations of the tunnel splittings with variations in energy or system parameters. The fluctuations were explained on the basis of a tunnel doublet, associated with regular regions in the phase space, invoved in an avoided crossing with a third state inhabiting the chaotic region of the phase space.

However, very recently classical nonlinear resonances in the near-integrable phase space regimes have also been implicated to play a dominant role in dynamical tunnelingrat1 ; rat2 . This has been dubbed, in analogy with CAT, as resonance-assisted tunneling (RAT) and studies revealed that tunnel splitting fluctuations, possibily more intense than in CAT, occur in near-integrable systems as well due to crossing of the regular tunnel doublets with a regular third state. The sheer richness of the phase space perspective in dynamical tunneling is further exemplified by the work of Frischat and Doronfrido . At present the consensus regarding dynamical tunneling seems to be that CAT and RAT are manifestations of the more general phenomenon of transport-assisted tunnelingfrido .

The preceeding discussion makes it clear that splitting fluctuations can occur in general nonintegrable systems and tunneling can be ascribed to chaos or resonances depending on the phase space nature of the third ’intruder’ state. It therefore seems natural to focus attention away from the avoided crossings in order to gain insights into the processrat2 . The present work is concerned with the study of dynamical tunneling in molecules from such a perspective. An essential difference from the work of Brodier, Schlagheck and Ullmorat2  is that the focus in this work is entirely on a perturbative evaluation of the splittings to quantitative accuracies. The splittings $`\mathrm{\Delta }`$ associated with the doublet states $`|i`$ and symmetry related $`|f`$ are calculated by high order nondegenerate perturbation theory involving the unperturbed states $`\{\alpha _r\}`$ as

$$\frac{\mathrm{\Delta }}{2}=\underset{ri,f}{}\frac{V_{i1}^{(1)}V_{12}^{(2)}\mathrm{}V_{nf}^{(n)}}{\mathrm{\Delta }E_{i1}\mathrm{\Delta }E_{i2}\mathrm{}\mathrm{\Delta }E_{nf}},$$
(1)
where $`\widehat{V}^{(r)}`$ are the local perturbations that connect the intermediate states through $`V_{qr}^{(r)}=\alpha _q|\widehat{V}^{(r)}|\alpha _r`$ with $`\mathrm{\Delta }E_{ir}=(E_iE_r)`$. The sum over the intermediate states can be viewed as a sum over paths in the state space and the number of intermediate states is related to the length of a path. The perturbative treatment is valid as long as $`V_{qr}^{(r)}/2<<\mathrm{\Delta }E_{ir}`$ and hence the method cannot be used to calculate splittings for states involved in avoided crossings. In principle, an infinite number of high order perturbative chains (paths in state space) exist which connect the two degenerate states. The hope has been that perhaps a few such paths would dominate the dynamical tunneling in the integrable and near-integrable cases. In this work we show that even in near-integrable cases there is no single dominant perturbative path. However, it is possible to identify a certain family of perturbative chains (“minimal paths” in state space) that are sufficient to provide quantitative results. Parameteric variations tune the phase space from the near-integrable limit to the mixed regime and it is observed that the minimal paths are not sufficient anymore to reproduce the splittings.

The purpose of this letter is to illustrate the observations above and accentuate the various mechanisms of dynamical tunneling. The system of choice is the water molecule which can be described by an effective spectroscopic Hamiltonianbagg

$$\widehat{H}_{\mathrm{eff}}=\widehat{H}_0+g\widehat{V}_{1:1}^{(12)}+\gamma \widehat{V}_{2:2}^{(12)}+\beta (\widehat{V}_{2:1}^{(1b)}+\widehat{V}_{2:1}^{(2b)}),$$
(2)
where $`\widehat{H}_0`$ is diagonal in the number basis $`(n_1,n_2,n_b)`$ and $`g,\gamma ,\beta `$ represent the strengths of the various perturbations. The various parameters of $`\widehat{H}_{\mathrm{eff}}`$ are given in the work of Baggottbagg . The zeroth order quantum numbers $`(n_1,n_2,n_b)`$ represent the excitation quanta in the two, equivalent OH-stretches and the bend mode respectively. The perturbations $`\widehat{V}`$ are off-diagonal in the number basis and have the form:

$$\widehat{V}_{m:n}^{(ij)}=[(\widehat{a}_i)^n(\widehat{a}_j^{})^m+\mathrm{h}.\mathrm{c}.]$$
(3)
where $`\widehat{a}_i`$ and $`\widehat{a}_i^{}`$ are the annhilation and creation operators for the mode $`i`$. The resonant perturbations $`\widehat{V}_{m:n}^{(ij)}`$ are responsible for the exchange of quanta (energy) between the modes $`i`$ and $`j`$. Classical-quantum correspondence studiesksjcp  have established that the classical limit Hamiltonian corresponding to $`\widehat{H}_{\mathrm{eff}}`$ is nonlinear and multiresonant. The location of the various resonance zonesksjcp  in state space are shown in Fig. (1). Note that the previous studieslawch ; stumar  on water correspond to integrable phase space regimes.

Clearly $`\widehat{H}_{\mathrm{eff}}`$ is symmetric under $`n_1n_2`$ and $`Pn_1+n_2+n_b/2`$ is a conserved quantity. Consequently $`\widehat{H}_{\mathrm{eff}}`$ for a given $`P`$ is of the order $`(P+1)(P+2)/2`$ making the numerical investigations rather straightforward. In addition $`\widehat{H}_{\mathrm{eff}}`$ exhibits a wide range of behaviour from integrable to near-integrable to mixed classical dynamics. Due to the exact $`C_{2v}`$ symmetry tunneling doublets exist over the entire range of the classical dynamics. In this work six tunneling doublets with $`P=8`$, denoted by $`|m,\pm m,n_b`$ with $`m=n_1+n_2`$ are studied. The labels used for the states are appropriate in the absence of the 2:1 resonances. We will continue to use them for convenience keeping in mind that the actual assignments of these states are differentksjcp . The energies of the states approximately span the range $`[24256,25574]`$ cm-1 above the ground state. For comparison the various resonant strengths have the typical values $`\beta 27`$ cm-1, $`g50`$ cm-1 and $`\gamma 1.0`$ cm-1.

Reperesentative phase space sections are shown in Fig. (2) for the highest energy state among the doublets considered. In Fig. (3) the doublet splittings, $`\mathrm{\Delta }_m`$, are shown for the various states $`|m,\pm m,n_b`$. Note that the states in Fig. (3) are not energy ordered and in particular the state $`|8,\pm 8,0`$ is lower in energy than the state $`|5,\pm 5,6`$. The splittings for the integrable cases (only 1:1 and 1:1+2:2) are shown in the inset to Fig. (3). As expectedstumar  the splittings show a monotonic decrease with increasing amount of stretch excitation and no fluctuations are observed. Even within the integrable cases it is relevant to note that the addition of the weak 2:2 results in a significant increase in $`\mathrm{\Delta }`$. High order perturbation calculationstumar  essentially reproduces the splittings to quantitative accuracy as is apparent from the inset in Fig. (3). Inclusion of the 2:1 stretch-bend resonances renders the system nonintegrable and clear deviations from the integrable cases are seen. In the main part of Fig. (3) the splittings are shown as one proceeds from near-integrable to the full, mixed scenario. The nature of the underlying phase space clearly reflects the transition from near-integrable to the mixed regime (cf. Fig. (2b-d)). The main difference from the integrable case is the nonmonotonic behaviour of $`\mathrm{\Delta }_m`$. The intense fluctuations for the state $`|7,\pm 7,2`$ are associated with its proximity to an avoided crossing. The key observation here is that the fluctuation due to avoided crossing is a robust feature and persists in the full system. Indeed it has been confirmed that in the 2:1 only case there is a strong doublet-doublet crossing which transforms into a singlet-doublet crossing in the full system and for the deuterated analog D2O a similar calculation shows absence of fluctuations in the full system despite the presence of significant chaostbp . Thus it is crucial to shift the focus away from the fluctuations to uncover the effect of chaos and resonances on dynamical tunneling in the system.

In Fig. (3) the phenomenon of RAT is evident in going from the 2:1 only system to the 2:1+2:2 system. The enhancements of the splittings are even more dramatic considering the fact that states with odd $`m`$ do not split at all if only the 2:2 resonance is present. Similarly the enhancement of the splittings on going from the 2:1+2:2 to the full system can be ascribed to CAT. Note that the enhancement due to CAT are typically an order of magnitude smaller than those due to RAT as has been earlier observed in the annular billiard studiesfrido . A crucial difference between CAT and RAT arises from the perturbative viewpoint. In order to highlight this difference we begin by noting that the structure of the resonances imply, in terms of order, $`\beta ^4g^2\gamma `$. Hence a consistent perturbative calculation of the splitting for the state $`|m,\pm m,n_b`$ is obtained as:

$$\mathrm{\Delta }_m=\underset{a,b,c}{}\beta ^ag^b\gamma ^c\underset{n}{}\mathrm{\Delta }_m(\mathrm{\Gamma }_{abc}^{(n)})$$
(4)
where $`n`$ indexes all possible paths $`\mathrm{\Gamma }_{abc}`$ for a particular choice of $`a,b,c`$ satisfying the constraint $`a+2b+4c=2m`$. It is clear that all possible paths have not been included in the above calculation of $`\mathrm{\Delta }_m`$. The paths included, called minimal paths, are all of the effective order $`2m`$. Examples of the minimal paths $`\mathrm{\Gamma }_{600}`$ and $`\mathrm{\Gamma }_{221}`$ relevant for calculating $`\mathrm{\Delta }_3`$ and $`\mathrm{\Delta }_5`$ respectively are shown in Fig. (1). One of the observations from this work is that for near-integrable systems a particular family of the minimal paths is dominant. In contrast, more than one family and higher order paths are required in the presence of significant chaos. As an example consider the case of $`\mathrm{\Delta }_5`$. This state is involved in a weak avoided crossing and hence provides a stringent test for the nondegenerate perturbation theory. In the near-integrable cases the family of paths $`\mathrm{\Gamma }_{1000}`$, $`\mathrm{\Gamma }_{202}`$, and $`\mathrm{\Gamma }_{601}`$ are sufficient. Perturbative calculations yield the splittings $`\mathrm{\Delta }_5(\mathrm{\Gamma }_{1000})=6.1\times 10^5`$ cm-1 and $`\mathrm{\Delta }_5(\mathrm{\Gamma }_{202,601})=7.2\times 10^3`$ cm-1 as compared to the numerically exact values of $`1.7\times 10^5`$ and $`6.2\times 10^3`$ cm-1 respectively. As mentioned earlier the 2:1 only system is less accurate due to the avoided crossing. In the integrable cases the perturbative results $`\mathrm{\Delta }_5(\mathrm{\Gamma }_{050})7.0\times 10^2`$ and $`\mathrm{\Delta }_5(\mathrm{\Gamma }_{050,031,012})0.23`$ cm-1 are in good agreement with the exact values of $`6.7\times 10^2`$ and $`0.23`$ cm-1 respectively. However, for the full system neither the near-integrable nor the integrable families are sufficient to reproduce the splitting ($`1.6`$ cm-1). In fact the families $`\mathrm{\Gamma }_{430},\mathrm{\Gamma }_{240},\mathrm{\Gamma }_{221}`$ play an important role and yield a splitting of about $`3.4`$ cm-1. For mixed systems, in general, a large number of families in the minimal set as well as higher order paths have to be included for a quantitatively accurate calculation of $`\mathrm{\Delta }_m`$. Thus, for instance, a minimal path calculation of $`\mathrm{\Delta }_6`$ yields a splitting which is about an order of magnitude large. The correlation of the preceeding observation with increasing amount of chaos in the phase space is an interesting possibilitytbp . As an example we mention that tuning the parameter $`\gamma `$ from about $`14`$ cm-1 turns the near-integrable $`2:1+2:2`$ system into a mixed system. The splitting $`\mathrm{\Delta }_5`$ is seen to increase by about a factor of $`22`$ whereas the minimal perturbative calculation predicts an accelaration by a factor of $`12`$. Clearly one anticipates higher order paths to play a significant role in mixed systems.

In conclusion we have shown that molecular systems with symmetry are prime candidates to study both chaos and resonance assisted dynamical tunnelingorti . High resolution spectroscopic studies of rovibrationally excited states are typically analysed in terms of such effective Hamiltonianshrs . The splitting patterns are crucial for understanding the extent of IVR in the molecule and hence the importance of dynamical tunnelingejhsar . This work is the first one to perform a detailed analysis of the possible mechanisms of dynamical tunneling in molecules with the hope that both CAT and RAT can be experimentally observed in the high resolution spectra. We have argued that a clear distinction between CAT and RAT can be established by focusing attention away from the avoided crossings. In addition high order perturbation theory, when valid, is quite successful in reproducing the tunnel splittings in the case of integrable and near-integrable systems. A single family of perturbative paths is sufficient but it is important to point out that among the many individual paths possible there is no single dominant path. For instance in the calculation of $`\mathrm{\Delta }_5`$ the family $`\mathrm{\Gamma }_{1000}`$ has $`242`$ paths which come with both positive and negative contributions and one needs to sum over all the paths to obtain accurate splittings. The indications are that more than one family of paths including those of order larger than the minimal order is needed as soon as there is a significant amount of chaos in the system. The situation is reminiscent of the multistep indirect paths invoked in the annular billiard studiesfrido . This provides important clues to the underlying mechanism of dynamical tunneling in the system. However, a semiclassical viewpoint on these perturbative minimal paths would lead to better insights into the process of dynamical tunneling from a phase space viewpoint.

It is a great pleasure to acknowledge Arul Lakshminarayan for helpful and critical discussions. This work was supported by funds from the Department of Science and Technology and the Council for Scientific and Industrial Research, India.