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# 1 A graph of 𝐹 versus 𝑘 is shown. For 0≤𝑘≤𝜋/2, the upper curve corresponds to the optimal quantum scheme and the lower curve represents the measurement scheme. For 𝜋/2≤𝑘≤𝜋, both measurement and optimal quantum schemes yield the identical results.

Universal Manipulation of a Single Qubit

Lucien Hardy and David D. Song

Centre for Quantum Computation

Clarendon Laboratory, Department of Physics

University of Oxford, Parks Road, Oxford OX1 3PU, U.K.

Abstract

We find the optimal universal way of manipulating a single qubit, $`|\psi (\theta ,\phi )`$, such that $`(\theta ,\phi )(\theta k,\phi l)`$. Such optimal transformations fall into two classes. For $`0k\pi /2`$ the optimal map is the identity and the fidelity varies monotonically from 1 (for $`k=0`$) to $`1/2`$ (for $`k=\pi /2`$). For $`\pi /2k\pi `$ the optimal map is the universal-NOT gate and the fidelity varies monotonically from $`1/2`$ (for $`k=\pi /2`$) to $`2/3`$ (for $`k=\pi `$). The fidelity $`2/3`$ is equal to the fidelity of measurement. It is therefore rather surprising that for some values of $`k`$ the fidelity is lower than $`2/3`$.

A unit of classical information is a bit, i.e. 0 or 1. Quantum information consists of qubits which are superposition of the states $`|0`$ and $`|1`$. Classical and quantum information differs in many ways. While classical information can be copied perfectly, the same is not true with qubits . Another feature that distinguishes classical and quantum information is a measurement. Unlike classical information, an unknown single qubit cannot be measured to give complete information about the qubit. In order to get a maximum information about an unknown qubit, we measure the qubit along any chosen basis $`\{|\varphi ,|\varphi ^{}\}`$. If the result is $`|\varphi `$, then we guess the unknown qubit to be $`|\varphi `$ and if the result is $`|\varphi ^{}`$, then we guess $`|\varphi ^{}`$. Averaging over all possible $`|\varphi `$’s (assuming a uniform distribution over the Bloch sphere), the fidelity is equal to 2/3. We cannot achieve a higher fidelity by using generalised measurements and hence 2/3 is the optimal measurement fidelity of an unknown qubit.

A qubit in Bloch vector notation is

$$|\psi (\theta ,\phi )\left(\begin{array}{c}\mathrm{cos}(\theta /2)\\ e^{i\phi }\mathrm{sin}(\theta /2)\end{array}\right)$$
(1)
The most general linear transformation on $`(\theta ,\phi )`$ is

$$(\theta ,\phi )(\theta k,\phi l)$$
(2)
with $`0k\pi `$ and $`0l2\pi `$. If $`\theta =0,\pi `$ then $`\phi `$ is undefined. For definiteness we will take $`\phi =0`$ in such cases. Since, when taking averages over the Bloch sphere, these anomalous cases are of measure zero we need not pay any special attention to them. This general transformation can be composed from two transformations. First

$$(\theta ,\phi )(\theta k,\phi )$$
(3)
and then $`\phi \phi l`$. The transformation on $`\phi `$ can be achieved perfectly by a unitary operation and so is of less interest to us. However, the transformation on $`\theta `$ cannot be achieved unitarily. To find the fidelity of general linear transformations of the form (2) it suffices to consider only the non-unitary part (3). We are interested only in universal transformations. These are those transformations for which the fidelity is independent of $`\theta `$ and $`\phi `$ of the input state. Furthermore, we will assume that the input distribution is uniform over the Bloch sphere. Since the area element $`\mathrm{sin}\theta d\theta d\phi `$ is not preserved in form (except when $`\theta =0,\pi `$) the output will not be uniform. In taking averages we integrate over a uniform distribution of the input variables which corresponds to integrating over a non-uniform distribution of the output variables.

Changing bits 0 to 1 and 1 to 0 is a NOT-gate in classical information case. In quantum case, changing $`|\psi =a|0+b|1`$ to $`|\psi ^{}=b^{}|0a^{}|1`$ requires anti-unitary transformation which is not allowed in quantum mechanics. In , it was shown that universal-NOT(U-NOT) operation can be achieved with 2/3 fidelity for a single input. This fidelity is same as the measurement fidelity. They showed that the U-NOT operation is no better than measuring a qubit first then preparing an orthogonal state. In Bloch vector notation the U-NOT gate corresponds to transforming $`|\psi (\theta ,\phi )`$ to $`|\psi (\theta \pi ,\phi )`$. This is a special case of the transformation (3) with $`k=\pi `$. Now consider the general case in which we transform $`|\psi (\theta ,\phi )`$ to $`|\psi (\theta k,\phi )`$. Naively, it may seem that the fidelity should be at least 2/3, since one could measure a qubit with $`2/3`$ fidelity and prepare a state in an appropriate direction. We will show in this paper that this is not so.

Let us take an example where $`k=3\pi /4`$. Therefore for a given unknown state $`|\psi `$, we want to prepare a state as close as possible to $`|\psi ^{}=|\psi (\theta 3\pi /4,\phi )`$. We choose a random state

$$|\varphi (u,v)\left(\begin{array}{c}\mathrm{cos}(u/2)\\ e^{iv}\mathrm{sin}(u/2)\end{array}\right)$$
(4)
and measure $`|\psi (\theta ,\phi )`$ in the basis of $`\{|\varphi ,|\varphi ^{}\}`$. If we get $`|\varphi `$ we prepare $`|\varphi ^{}|\varphi (u3\pi /4,v)`$ and if we get $`|\varphi ^{}`$, then $`|\varphi ^{}`$ is prepared. As a density matrix, the state we prepare by this method is

$$\rho ^{(1)}=|\psi |\varphi |^2|\varphi ^{}\varphi ^{}|+|\psi |\varphi ^{}|^2|\varphi ^{}\varphi ^{}|$$
(5)
We take the average of $`\rho ^{(1)}`$ over uniform distributions of $`|\varphi `$ on the Bloch sphere to obtain $`\overline{\rho ^{(1)}}`$ and then the fidelity is given by

$$F=\frac{1}{4\pi }_0^\pi _0^{2\pi }\psi ^{}|\overline{\rho ^{(1)}}|\psi ^{}\mathrm{sin}\theta d\theta d\phi =.5833\mathrm{}$$
(6)
This value is lower than 2/3. Can we do any better? If we prepare $`|\varphi ^{}`$ when the result is $`|\varphi `$ and prepare $`|\varphi `$ for $`|\varphi ^{}`$, i.e.

$$\rho ^{(2)}=|\psi |\varphi |^2|\varphi ^{}\varphi ^{}|+|\psi |\varphi ^{}|^2|\varphi \varphi |$$
(7)
then

$$F=\frac{1}{4\pi }_0^\pi _0^{2\pi }\psi ^{}|\overline{\rho ^{(2)}}|\psi ^{}\mathrm{sin}\theta d\theta d\phi =.6178\mathrm{}$$
(8)
This fidelity is still lower than 2/3 but higher than the value in (6). This rather surprising result is due to the different phase angles of $`|\psi `$ and $`|\varphi `$. Suppose $`|\psi |\varphi |^2=2/3`$, then the rotation $`|\psi (\theta ,\phi )|\psi (\theta \pi ,\phi )`$ and $`|\varphi (u,v)|\varphi (u\pi ,v`$ yields the same fidelity $`|\psi (\theta \pi ,\phi )|\varphi (u\pi ,v)|^2=2/3`$. However if the rotation is some other angle $`k\pi `$ or $`0`$, then $`|\psi (\theta k,\phi )|\varphi (uk,v)|^2`$ may not be 2/3 because $`\phi `$ and $`v`$ are not necessarily the same. If the phase angles $`\phi `$ and $`v`$ are the same, then for any $`k`$, $`|\psi (\theta k,\phi )|\varphi (uk,\phi )|^2=2/3`$. For $`\pi /2k\pi `$ in $`|\psi (\theta k,\phi )`$, $`\overline{\rho ^{(2)}}`$ yields the fidelity

$$F^{(2)}=\frac{1}{12}(6+\mathrm{cos}(\pi k)+\mathrm{cos}(\pi +k))$$
(9)
For $`0k\pi /2`$, we consider the usual measurement density matrix,

$$\rho ^{(3)}=|\psi |\varphi |^2|\varphi \varphi |+|\psi |\varphi ^{}|^2|\varphi ^{}\varphi ^{}|$$
(10)
and the fidelity is given as

$$F^{(3)}=\frac{1}{2}+\frac{1}{6}\mathrm{cos}k$$
(11)
For $`\pi /2k\pi `$, we will show that $`F^{(2)}`$ in (9) is indeed the optimal fidelity. $`|\psi (\theta k,\phi )`$ for $`0k\pi /2`$ can be obtained with better fidelity than $`F^{(3)}`$ in (11). We expect this since for $`k=0`$, the identity operation gives $`|\psi `$ with fidelity 1.

By considering the most general type of transformation on a single qubit we will find the one which maximises the fidelity for the transformation (3). We will follow the method of Buzek et al. . The most general operation available to us is to perform unitary evolution on the single qubit and some ancilla prepared in a known state $`|Q`$ (this is taken to be normalised). This gives

$`|0|Q`$ $``$ $`|1|A+|0|B`$
$`|1|Q`$ $``$ $`|0|\stackrel{~}{A}+|1|\stackrel{~}{B}`$ (12)
where $`|A,|\stackrel{~}{A},|B,|\stackrel{~}{B}`$ may not be normalised. From the normalisation and the orthogonality of (12), $`|A|^2+|B|^2=|\stackrel{~}{A}|^2+|\stackrel{~}{B}|^2=1`$ and $`A|\stackrel{~}{B}+B|\stackrel{~}{A}=0`$. We let $`|\psi `$ transform under (12), trace over the ancilla and obtain a density matrix $`\rho ^{(\mathrm{o}ut)}`$ for our qubit. The fidelity is then given by $`F=\psi ^{}|\rho ^{(\mathrm{o}ut)}|\psi ^{}`$ where $`|\psi ^{}=|\psi (\theta k,\phi )`$. When expressed explicitly in terms of $`\theta `$, $`\phi `$, $`|A`$, $`|B`$, $`|\stackrel{~}{A}`$, and $`|\stackrel{~}{B}`$ this expression for the fidelity has 48 terms. We can compare coefficients of those terms dependent on $`e^{\pm i\phi }`$ and $`e^{\pm 2i\phi }`$. These coefficients must vanish in order for the transformation in (12) to be independent of $`\phi `$ (which we require for universality), $`A|\stackrel{~}{A}=A|B`$ = $`\stackrel{~}{B}|\stackrel{~}{A}=\stackrel{~}{A}|B`$=$`\stackrel{~}{B}|A=0`$. Of those terms remaining, some have a dependence on $`\theta `$. These terms must also vanish (by universality). This leaves only two terms giving us an expression for the fidelity:

$$F=\mathrm{sin}^2\frac{k}{2}|\stackrel{~}{A}|^2+\mathrm{cos}^2\frac{k}{2}|\stackrel{~}{B}|^2=(\mathrm{cos}^2\frac{k}{2}\mathrm{sin}^2\frac{k}{2})|\stackrel{~}{B}|^2+\mathrm{sin}^2\frac{k}{2}$$
(13)
By comparing coefficients of functions of $`\theta `$ of those terms having a $`\theta `$ dependence and setting them to zero we obtain the following conditions

$`2|\stackrel{~}{A}|^22|\stackrel{~}{B}|^2+\stackrel{~}{B}|B+B|\stackrel{~}{B}=0`$ (14)
$`2|B|^22|A|^2\stackrel{~}{B}|BB|\stackrel{~}{B}=0`$ (15)
$`|A|^2+|\stackrel{~}{A}|^2|B|^2|\stackrel{~}{B}|^2+B|\stackrel{~}{B}+\stackrel{~}{B}|B=0`$ (16)
$`(\mathrm{cos}^2{\displaystyle \frac{k}{2}}2\mathrm{sin}^2{\displaystyle \frac{k}{2}})|\stackrel{~}{A}|^2+(\mathrm{sin}^2{\displaystyle \frac{k}{2}}2\mathrm{cos}^2{\displaystyle \frac{k}{2}})|\stackrel{~}{B}|^2+\mathrm{sin}^2{\displaystyle \frac{k}{2}}|B|^2+\mathrm{cos}^2{\displaystyle \frac{k}{2}}|A|^2`$
$`+(\mathrm{cos}^2{\displaystyle \frac{k}{2}}\mathrm{sin}^2{\displaystyle \frac{k}{2}})\stackrel{~}{B}|B+(\mathrm{cos}^2{\displaystyle \frac{k}{2}}\mathrm{sin}^2{\displaystyle \frac{k}{2}})B|\stackrel{~}{B}=0`$ (17)
From (14) and (15),

$$|A|^2=|\stackrel{~}{A}|^2,|B|^2=|\stackrel{~}{B}|^2$$
(18)
which then implies (17) is equal to (15). From (16), with $`\eta =\mathrm{R}e(B|\stackrel{~}{B})/|\stackrel{~}{B}|^2`$ (therefore $`|\eta |1`$),

$$|\stackrel{~}{B}|^2=\frac{1}{2\eta }$$
(19)
For $`\pi /2k\pi `$, $`|\stackrel{~}{B}|^2`$ needs to be minimum to give a maximum fidelity in (13). Therefore with $`\eta =1`$, the fidelity is

$$F=\frac{1}{3}\mathrm{cos}^2\frac{k}{2}+\frac{2}{3}\mathrm{sin}^2\frac{k}{2}$$
(20)
which is same as (9). Therefore for $`\pi /2k\pi `$, the measurement based preparation as in $`\rho ^{(2)}`$ is indeed optimal. The transformation satisfying (14-18) and (20) is same as the U-NOT transformation of Buzek et al. in . The fidelity in (20) has the highest value of 2/3 when $`k=\pi `$ and the lowest 1/2 when $`k=\pi /2`$. The graph for $`k`$ and $`F`$ is shown in Figure 1 where the measurement and the quantum schemes yield the identical result. For $`0k\pi /2`$, $`|\stackrel{~}{B}|^2`$ needs to be maximum to have a maximum fidelity in (13). Therefore with the choice of $`\eta =1`$,

$$F=\mathrm{cos}^2\frac{k}{2}$$
(21)
This transformation is simply a trivial identity map and it has maximum fidelity of 1 when $`k=0`$ and minimum 1/2 when $`k=\pi /2`$. A graph of $`k`$ and $`F`$ for $`0k\pi /2`$ is shown in Figure 1. In this case, the quantum scheme has higher fidelity than the measurement scheme.

It follows that for a general transformation linear in the spherical coordinates, namely $`(\theta ,\phi )(\theta k,\phi l)`$, the procedures which optimise fidelity fall into two distinct classes. (A) For $`0k\pi /2`$ the optimal procedure is the identity map which performs better than a measurement based scheme. In this range the maximum fidelity (equal to one) is achieved, not surprisingly, when $`k=0`$. (B) For $`\pi /2k\pi `$ the U-NOT transformation is optimal. This procedure performs only as well as a measurement based scheme. In this range the maximum fidelity (equal to $`2/3`$) is achieved, perhaps a little surprisingly, only for the case $`k=\pi `$ which, if $`l=0`$, corresponds to a universal NOT operation. Since $`\phi `$ can be varied linearly by a unitary transformation, $`l`$ can take any value in either of these two classes.

Acknowledgements

We are grateful to Leah Henderson and Vlatko Vedral for discussions on this topic. LH is funded by a Royal Society University Research Fellowship.


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# 1 Introduction

## 1 Introduction

Further progress in clarifying the relation between strings and large $`N`$ (non) supersymmetric gauge theories depends on better understanding of superstrings in $`AdS`$-type spaces with Ramond-Ramond (RR) fluxes . The simplest most symmetric example is provided by $`AdS_5\times S^5`$ background of type IIB superstring theory. One could hope that this theory, though quite nontrivial being dual to large N $`𝒩=4`$ Super Yang-Mills theory , should be, like the flat superstring, explicitly solvable.

To be able to treat this theory in an efficient way one is to use the Green-Schwarz (GS) approach  where one views the string as moving in a superspace $`(x^m,\theta _\alpha ^I)`$ with $`\theta `$’s being space-time spinors. Then the coupling of the string to RR background has a local form $`\overline{\theta }\mathrm{\Gamma }^{\mathrm{}}\theta xxF_{\mathrm{}}`$, just like its coupling to the curvature. It is useful to note that while the string “feels” the metric already at the classical level, it interacts with the RR background only through its quantum fluctuations.

There seems to be no alternative to the GS-type description: even the simplest one-loop computation in GS formalism like the one described in  would require a summation of infinite number of diagrams with any number of flat-space RR vertex operator insertions if done in NSR formalism. In fact, the use of GS action is very natural from “solitonic” point of view. If one views fundamental strings as “electric”-type objects appearing in $`D=10`$ supergravity, the effective action for collective coordinates of a long such string will naturally have a GS form, as dictated by the residual space-time (super)symmetries (the characteristic non-linear $`x\overline{\theta }\mathrm{\Gamma }^{\mathrm{}}\theta `$ kinetic term originates from the standard Dirac $`D=10`$ fermion action reduced to the world surface of the string, with Dirac matrices becoming the projected ones). The theory should of course contain not only long but also short strings which should also be described by the same manifestly space-time supersymmetric action.

A simple way to construct the flat space GS action is to view strings as propagating in flat coset superspace = (10-d super Poincare group)/( 10-d Lorentz group). The action is then a kind of Wess-Zumino-Witten action which has global supersymmetry and local $`\kappa `$-symmetry which ensures the correct number of the fermionic degrees of freedom. It is defined in terms of the basic objects which are the left-invariant Cartan 1-forms on the type IIB coset superspace $`G^1dG=L^{\widehat{A}}P_{\widehat{A}}+L^IQ_I,`$ $`L^{\widehat{A}}=dX^ML_M^{\widehat{A}},`$ $`X^M=(x,\theta ),`$ where $`[P_{\widehat{A}},P_{\widehat{B}}]=0,`$ $`\{Q_I,Q_J\}=2i\delta _{IJ}(𝒞\mathrm{\Gamma }^{\widehat{A}})P_{\widehat{A}}`$ and $`G=G(x,\theta )`$ is a coset representative $`G(x,\theta )=\mathrm{exp}(x^{\widehat{A}}P_{\widehat{A}}+\theta ^IQ_I)`$. The coset space vielbeins are given by $`L^{\widehat{A}}=dx^{\widehat{A}}i\overline{\theta }^I\mathrm{\Gamma }^{\widehat{A}}d\theta ^I,`$ $`L^I=d\theta ^I.`$ Here $`x^{\widehat{A}}`$ are flat bosonic coordinates and $`\theta ^I`$ ($`I=1,2`$) are two left Majorana-Weyl 10-d spinors. The general form of the action is

$$S=\frac{1}{2}_{M_3}d^2\sigma \sqrt{g}g^{\mu \nu }L_\mu ^{\widehat{A}}L_\nu ^{\widehat{A}}+i_{M_3},=s^{IJ}L^{\widehat{A}}(\overline{L}^I\mathrm{\Gamma }^{\widehat{A}}L^J),$$
(1)
where $`s^{IJ}\mathrm{diag}(1,1)`$, $`2\pi \alpha ^{}=1`$ and $`d=0`$. The 2-d metric $`g_{\mu \nu }`$ ($`\mu ,\nu =0,1`$) has signature $`(,+)`$, and $`g\text{det}g_{\mu \nu }`$. The first “kinetic” term in the action has “degenerate” metric, containing only the square of translational but not spinor Cartan forms, as the latter would lead to an action quadratic in fermionic derivatives and thus to potential non-unitarity. The explicit 2-d component form of the GS action is

$$S_0=d^2\sigma [\frac{1}{2}\sqrt{g}g^{\mu \nu }(_\mu x^{\widehat{A}}i\overline{\theta }^I\mathrm{\Gamma }^{\widehat{A}}_\mu \theta ^I)(_\nu x^{\widehat{A}}i\overline{\theta }^J\mathrm{\Gamma }^{\widehat{A}}_\nu \theta ^J)$$

$$iϵ^{\mu \nu }s^{IJ}\overline{\theta }^I\mathrm{\Gamma }^{\widehat{A}}_\nu \theta ^J(_\mu x^{\widehat{A}}\frac{1}{2}i\overline{\theta }^K\mathrm{\Gamma }^{\widehat{A}}_\mu \theta ^K)].$$
(2)
The key observation that allows one to construct the GS action for superstrings in $`AdS_5\times S^5`$ is that this bosonic space is a coset

$`G/H=SO(2,4)/SO(1,4)\times SO(6)/SO(5)`$

and also has maximal supersymmetry. As a result, one can consider strings moving not in flat superspace but in supercoset space with the right bosonic part and right number supersymmetries: one is to replace the 10-d super Poincare group by a smaller one $`PSU(2,2|4)`$ and the Lorentz group $`SO(1,9)`$ by its subgroup $`SO(1,4)\times SO(5)`$, i.e. to consider the supercoset $`PSU(2,2|4)/[SO(1,4)\times SO(5)]`$ , which still has 10 bosonic and 32 Grassmann dimensions.

The bosonic part of the GS action is simply the standard symmetric space sigma model $`L=\mathrm{Tr}[(g^1g)_{G/H}]^2`$. This is obviously not a conformal sigma model; it is the addition of the fermions $`\theta `$ that converts it into a conformal 2-d model of a novel type. The fermions couple the $`AdS_5`$ and $`S^5`$ parts of the action together and contribute the RR 5-form term to the beta-function of the metric making it vanish , $`R_{mn}(F_5F_5)_{mn}=0`$.

The superalgebra $`psu(2,2|4)`$ plays the central role in the construction of the GS action in $`AdS_5\times S^5`$ . Its even part is the sum of the algebra $`so(4,2)`$ which is the isometry algebra of $`AdS_5`$ and the algebra $`so(6)`$ which is the isometry algebra of $`S^5`$. The odd part consists of 32 supercharges corresponding to 32 Killing spinors in $`AdS_5\times S^5`$ vacuum  of type IIB supergravity . In the “5+5” split its generators are: $`(P_A,J_{AB};P_A^{},J_{A^{}B^{}};Q_{\alpha \alpha ^{}I})`$, i.e. the the two sets of translations and rotations in $`so(2,4)`$ and $`so(6)`$ ($`A=0,1,\mathrm{},4;A^{}=1,\mathrm{}5`$) and 32 supercharges ($`\alpha =1,2,3,4;\alpha ^{}=1,2,3,4`$). The commutation relations of $`psu(2,2|4)`$ in $`so(4,1)so(5)`$ basis are  $`[P_A,P_B]=J_{AB},`$ $`[P_A^{},P_B^{}]=J_{A^{}B^{}},`$ $`[J^{AB},J^{CE}]=\eta ^{BC}J^{AE}+\mathrm{},`$ $`[J^{A^{}B^{}},J^{C^{}E^{}}]=\eta ^{B^{}C^{}}J^{A^{}E^{}}+\mathrm{},`$ $`[Q_I,P_A]=\frac{i}{2}ϵ_{_{IJ}}Q_J\gamma _A,`$ $`[Q_I,J_{AB}]=\frac{1}{2}Q_I\gamma _{AB},`$ $`[Q_I,P_A^{}]=\frac{1}{2}ϵ_{_{IJ}}Q_J\gamma _A^{},`$ $`[Q_I,J_{A^{}B^{}}]=\frac{1}{2}Q_I\gamma _{A^{}B^{}},`$ and

$$\{Q_{\alpha \alpha ^{}I},Q_{\beta \beta ^{}J}\}=\delta _{_{IJ}}[2iC_{\alpha ^{}\beta ^{}}(C\gamma ^A)_{\alpha \beta }P_A+2C_{\alpha \beta }(C^{}\gamma ^A^{})_{\alpha ^{}\beta ^{}}P_A^{}]$$

$$+ϵ_{_{IJ}}[C_{\alpha ^{}\beta ^{}}(C\gamma ^{AB})_{\alpha \beta }J_{AB}C_{\alpha \beta }(C^{}\gamma ^{A^{}B^{}})_{\alpha ^{}\beta ^{}}J_{A^{}B^{}}],$$
(3)
where $`C,C^{}`$ are charge conjugation matrices. Then $`G^1dG=L^AP^A+\frac{1}{2}\widehat{L}^{AB}\widehat{J}^{AB}`$ $`+L^A^{}P^A^{}+\frac{1}{2}L^{A^{}B^{}}J^{A^{}B^{}}+L^{I\alpha \alpha ^{}}Q_{I\alpha \alpha ^{}},`$ and the coset representative can be parametrised as $`G(x,\theta )=g(x)\mathrm{exp}(\theta Q)`$.

Following the same steps as above in the construction of the flat space GS action one finds  that there exists a unique action of the type (1) which has the required properties: its bosonic part is sigma model on $`AdS_5\times S^5`$ , it has local $`\kappa `$-symmetry and global $`PSU(2,2|4)`$ symmetry, and it reduces to the standard GS action in the flat space limit. The action has the same structure as (1) but now in terms of the 1-forms corresponding to the supercoset $`PSU(2,2|4)/[SO(1,4)\times SO(5)]`$. The leading terms in the resulting action are essentially the covariantisation of the flat space action with the ordinary derivatives on $`\theta `$’s replaced by generalized $`AdS_5\times S^5`$ covariant ones (containing extra terms linear in Dirac matrices that correspond to the interaction with RR background).

As in the standard WZW model case it is possible to argue that this action defines a conformal sigma model : (i) the WZ term in the action is not renormalized since the divergences are local and manifestly covariant, while the WZ term is not; (ii) the kinetic $`L^2`$ term in the action can be renormalized only by an overall constant because of the global symmetries of the sigma model; (iii) assuming that $`\kappa `$-symmetry is preserved by regularization, it should relate the coefficients of the kinetic andthe WZ terms as it does at the classical level (i.e. $`\kappa `$-symmetry plays the role of the affine symmetry in WZW model).

The explicit form of all higher order terms in $`\theta `$ in the action of  can be found in . Fixing a “covariant” $`\kappa `$-symmetry gauge (e.g. $`\theta ^1=\mathrm{\Gamma }_{0123}\theta ^2`$) one gets a relatively simple action which contains terms of quadratic and quartic order in fermions only . The resulting action has the same feature as the flat space GS action: its fermionic kinetic term is coupled to derivative of all bosonic string coordinates. It is non-degenerate when expanded near “long” string configuration , and thus the GS action provides a well-defined and useful tool for computing quantum corrections to “long” string configurations ending on Wilson loops  at the boundary of $`AdS_5`$ .

The bosonic part of the action for a string in $`AdS_5\times S^5`$ is the sigma model corresponding to the space-time metric ($`m=1,\mathrm{},4`$)

$$ds^2=R^2\left[\frac{1}{w^2}\left(dw^2+dx^mdx^m\right)+d\mathrm{\Omega }_5^2\right].$$
(4)
The $`\alpha ^{}`$ expansion corresponds to the large ’t Hooft coupling $`\lambda `$ expansion in the dual gauge theory, with the leading classical contribution proportional to $`\frac{R^2}{\alpha ^{}}=\sqrt{\lambda }`$ .

## 2 Quadratic approximation

The part of the action in  quadratic in $`\theta ^I`$ is a direct generalization of the quadratic term in the flat-space GS action (2)

$$S_F^{(2)}=id^2\sigma (\sqrt{g}g^{\mu \nu }\delta ^{IJ}ϵ^{\mu \nu }s^{IJ})\overline{\theta }^I\rho _\mu D_\nu \theta ^J.$$
(5)
Here $`\rho _\mu `$ are projections of the 10-d Dirac matrices,

$$\rho _\mu \mathrm{\Gamma }_{\widehat{m}}E_M^{\widehat{m}}_\mu x^M=(\mathrm{\Gamma }_AE_M^A+\mathrm{\Gamma }_A^{}E_M^A^{})_\mu x^M,$$
(6)
and $`E_M^{\widehat{m}}`$ is the vielbein of the 10-d target space metric. The covariant derivative $`D_\mu `$ is the projection of the 10-d derivative $`D_M=_M+\frac{1}{4}\omega _M^{\widehat{m}\widehat{n}}\mathrm{\Gamma }_{\widehat{m}\widehat{n}}\frac{1}{85!}\mathrm{\Gamma }^{\widehat{m}_1\mathrm{}\widehat{m}_5}\mathrm{\Gamma }_Me^\mathrm{\Phi }F_{\widehat{m}_1\mathrm{}\widehat{m}_5}`$ which appears in the Killing spinor equation of type IIB supergravity,

$$D_\mu \theta ^I\left(\delta ^{IJ}\mathrm{D}_\mu \frac{i}{2}ϵ^{IJ}\stackrel{~}{\rho }_\mu \right)\theta ^J,\stackrel{~}{\rho }_\mu \left(\mathrm{\Gamma }_AE_M^A+i\mathrm{\Gamma }_A^{}E_M^A^{}\right)_\mu x^M,$$
(7)
where $`\mathrm{D}_\mu =_\mu +\frac{1}{4}_\mu x^M\omega _M^{\widehat{m}\widehat{n}}\mathrm{\Gamma }_{\widehat{m}\widehat{n}}`$ and the term with $`\stackrel{~}{\rho }_\mu `$ originates from the coupling to the RR 5-form field strength (note that $`F_5`$ is proportional to the $`ϵ`$-tensors on $`AdS_5\times S^5`$ ). The presence of the generalized “Killing spinor” covariant derivative explains why the action has 32 global supersymmetries. Fixing the natural type IIB $`\kappa `$-symmetry gauge $`\theta ^1=\theta ^2\theta `$ one finds from (5)

$$S_F^{(2)}=2id^2\sigma \left(\sqrt{g}g^{\mu \nu }\overline{\theta }\rho _\mu \mathrm{D}_\nu \theta \frac{i}{2}ϵ^{\mu \nu }\overline{\theta }\rho _\mu \stackrel{~}{\rho }_\nu \theta \right).$$
(8)
This action gives a well-defined fermion kinetic term if one expands near generic “long” string configuration. It is easy to check directly that the bosonic $`AdS_5\times S^5`$ sigma model becomes indeed 1-loop finite when supplemented by the fermionic action (8)  with the “mass term” in (8) originating from the RR coupling playing the crucial role.

The kinetic term in GS action in curved target space background (8) involving $`\theta `$’s which are world-volume scalars can be transformed into the “2-d spinor” form, i.e. into the standard kinetic term for a set of 2-d Dirac fermions defined on a curved 2-d space. In the flat space case expanding the GS action near $`x^0=\tau ,x^1=\sigma `$ one concludes that (e.g. in the $`\theta ^1=\theta ^2`$ gauge) the kinetic fermionic term takes the form $`\overline{\theta }\mathrm{\Gamma }^a_a\theta `$, $`a=0,1`$, so that choosing the representation of 10-d Dirac matrices where $`\mathrm{\Gamma }^{0,1}`$ are proportional to 2-d Dirac matrices one may re-interpret the action as the one for a collection of 2-d Dirac fermions.

In general, to achieve a similar transformation one is to apply a local target space Lorentz rotation to GS $`\theta `$ (as discussed previously mostly in flat space in ). Identifying the 2-d metric $`g_{\mu \nu }`$ with the induced metric $`G_{mn}_\mu x^m_\nu x^n`$ we can write the first covariant derivative term in (8) as $`id^2\sigma \sqrt{g}g^{\mu \nu }`$ $`(\overline{\theta }\rho _\mu _\nu \theta _\mu \overline{\theta }\rho _\nu \theta ).`$ Introducing the tangent $`t_\alpha ^m`$ ($`m=0,1,\mathrm{},9`$, $`\alpha =0,1`$) and normal $`n_s^m`$ ($`s=1,\mathrm{},8`$) vectors to the world surface which form orthonormal 10-d basis ($`g_{\mu \nu }=e_\mu ^\alpha e_\nu ^\beta \eta _{\alpha \beta }`$), i.e. $`t_\alpha ^m=e_\alpha ^\mu _\mu \overline{x}^m,`$ $`(t_\alpha ,t_\beta )=\eta _{\alpha \beta },`$ $`(t_\alpha ,n_s)=0,`$ $`(n_s,n_u)=\delta _{su},`$ where $`(a,b)=G_{mn}a^mb^n`$, one can make a local $`SO(1,9)`$ rotation of this basis which transforms the set of $`\sigma `$-dependent 10-d Dirac matrices into the 10 constant Dirac matrices $`\rho _\alpha (\sigma )=e_\alpha ^\mu \rho _\mu =S(\sigma )\mathrm{\Gamma }_\alpha S^1(\sigma ),`$ $`\rho _s(\sigma )=n_s^mE_m^{\widehat{a}}\mathrm{\Gamma }_{\widehat{a}}=S(\sigma )\mathrm{\Gamma }_sS^1(\sigma ).`$ One may further choose a representation in which $`\mathrm{\Gamma }_\alpha =\tau _\alpha \times I_8`$, where $`\tau _\alpha `$ are 2-d Dirac matrices. Depending on specific embedding and particular curved target space metric, one may then be able to write the action (5) as the action for 2-d Dirac fermions coupled to curved induced 2-d metric and interacting with some 2-d gauge fields (coming from $`S^1dS`$).

Simple examples of when this happens  will be discussed below. We shall consider embeddings of the string world sheet into the $`AdS_3`$ part of the $`AdS_5`$ space, so that there will be only one normal direction and the extra normal bundle 2-d gauge connection will be absent.

## 3 One-loop partition function for “straight” string

The simplest classical solution for string in $`AdS_5\times S^5`$ is a straight string with the world surface spanned by the radial direction of $`AdS_5`$ and time. The Euclidean solution and the corresponding induced metric are

$$x^0=\tau ,x^4w=\sigma ,ds^2=\frac{1}{\sigma ^2}(d\tau ^2+d\sigma ^2).$$
(9)
The induced metric on the world sheet is that of $`AdS_2`$, with constant negative curvature $`R^{(2)}=2`$ (we set the radius $`R`$ of $`AdS_5`$ to $`1`$). This solution corresponds to a single straight Wilson line at the boundary running along the Euclidean time direction. This is a BPS object in string theory (and boundary gauge theory): it corresponds to a static fundamental string stretched between a single D3-brane (placed at the boundary of $`AdS_5`$) and $`N`$ coinciding D3-branes (placed at the horizon and supporting $`AdS_5\times S^5`$ by their RR flux). Therefore, the partition function should be equal to 1. The properly defined (subtracted) classical string action evaluated on this background vanishes, and the corresponding 1-loop correction to the vacuum energy defined with respect to a certain time-like Killing vector vanishes too (see below). Relating the vacuum energy to the partition function using a conformal rescaling argument one concludes that $`Z=1`$ .111It should be mentioned that while the properly defined vacuum energy of a supersymmetric field theory in $`AdS`$ space should vanish, this does not automatically imply (in contrast to what happens in flat space) that the partition function of such theory should be equal to 1 (cf. ). The calculation of the partition function is rather subtle, and depends on a regularization prescription and proper inclusion of measure factors. For practical applications, the precise value of $`Z`$ (which is simply a constant) is not actually important, and one may normalize with respect to it in computing $`Z`$ for more general string configurations. Any smooth Wilson loop looks in the UV region like a straight line. This translates into the behavior of the minimal surface near the boundary of $`AdS_5`$ space. In the general case one will have to calculate the partition function for a complicated two dimensional field theory. But asymptotically the minimal surface will approach $`AdS_2`$, and the small fluctuation operators (in particular, the asymptotic values of the masses of the fluctuation fields) will also be the same as for a straight string. Subtleties related to divergences and asymptotic boundary conditions can be automatically avoided in more general cases by normalizing with respect to the partition function of the straight string.

The bosonic part of the action for small fluctuations in conformal gauge turns out to be ($`a,b`$ are the tangent indices of $`AdS_5`$ part and $`p`$ of $`S^5`$ part)

$$S_{2\mathrm{B}}=\frac{1}{2}d^2\sigma \sqrt{g}\left(D^\mu \zeta ^aD_\mu \zeta ^a+X_{ab}\zeta ^a\zeta ^b+D^\mu \zeta ^pD_\mu \zeta ^p\right),$$
(10)
where $`X_{ab}=\mathrm{diag}(1,2,2,2,1)`$. The only non-trivial components of the covariant derivative are $`D_0\zeta ^0=_0\zeta ^0w^1\zeta ^4,D_0\zeta ^4=_0\zeta ^4+w^1\zeta ^0.`$ Because of the direct embedding of the world sheet into the target space the projection of the target space connection on the world sheet is the same as the spin connection of the induced metric and thus the action of the conformal gauge ghosts is identical to the action of the longitudinal modes $`\zeta ^0`$, $`\zeta ^4`$.

The quadratic part of the fermion action in the $`\theta ^1=\theta ^2`$ gauge (8) is

$$L_{2\mathrm{F}}=2i\sqrt{g}\left(\overline{\theta }\rho ^\mu \widehat{}_\mu \theta +i\overline{\theta }\rho _3\theta \right),\rho _3\frac{1}{2}ϵ^{\alpha \beta }\rho _\alpha \rho _\beta =\mathrm{\Gamma }_{04},$$
(11)
where $`\rho _\alpha =(\mathrm{\Gamma }_0,\mathrm{\Gamma }_4)`$ may be identified with (2-d Dirac matrices)$`\times I_8`$. Thus the quadratic fermionic part of GS action has exactly the same form as the action for 2-d fermions in curved $`AdS_2`$ space with a mass term. Assuming the standard $`d^2\sigma \sqrt{g}\overline{\theta }\theta `$ normalization, the corresponding Dirac operator is $`i\rho ^\mu \widehat{}_\mu \rho _3=iw(\mathrm{\Gamma }_0_0+\mathrm{\Gamma }_4_1)\frac{1}{2}i\mathrm{\Gamma }_4\mathrm{\Gamma }_0\mathrm{\Gamma }_4`$ and its square is $`\widehat{}^2+\frac{1}{4}R^{(2)}+1`$.

Ignoring the ghosts and longitudinal modes222 Their contributions cancel each other modulo a constant related to difference in boundary conditions which is crucial for making the full string partition function finite . we are left with a 2-d field theory on $`AdS_2`$ containing 5 massless scalars, 3 scalars with (mass)2 =2, and eight fermions with (mass)2 =1. Supersymmetric field theories on $`AdS_2`$ were studied, e.g., in . The fields in the $`𝒩=1`$ scalar supermultiplet in $`AdS_2`$ may have the following bosonic and fermionic masses : $`m_B^2=\mu ^2\mu ,m_F=\mu ,`$ ($`\mu `$ is a free parameter). We thus have 5 “massless” multiplets with $`\mu =1`$ ($`m_B^2=0`$, $`m_F=1`$) and 3 “massive” multiplets with $`\mu =1`$ ($`m_B^2=2`$, $`m_F=1`$). It is possible to combine a $`\mu =1`$ and a $`\mu =1`$ multiplet into an $`𝒩=2`$ multiplet (the dimensional reduction of the 4-d chiral multiplet to 2 dimensions). Two $`\mu =1`$ multiplets also form an $`𝒩=2`$ multiplet (a dimensional reduction of the 4-d vector multiplet). Three chiral and one vector multiplets in $`D=4`$ make one $`𝒩=4`$ vector in four dimensions, and that suggests that the 8 scalars and 8 fermions that we have should form one $`𝒩=8`$ multiplet in two dimensions (cf. ). We finally obtain the following partition function for the transverse and fermionic modes with the scalar and spinor Laplace operators defined with respect to the Euclidean $`AdS_2`$ metric with radius 1 ($`R^{(2)}=2`$)

$$Z_{B+F}=\frac{\text{det}^{8/2}\left(\widehat{}^2+\frac{1}{4}R^{(2)}+1\right)}{\text{det}^{3/2}\left(^2+2\right)\text{det}^{5/2}\left(^2\right)}.$$
(12)
One should impose proper boundary conditions consistent with $`AdS_n`$ supersymmetry . Those imply that the resulting spectra of Laplace operators are discrete in spatial direction (and not continuous as one would expect in a non-compact hyperbolic space).

Instead of calculating the partition function directly, we may start with the vacuum energy. It is given by the determinant of the operator scaled to remove the factor of $`g^{00}`$ from in front of $`_t^2`$ (see ). Changing the world sheet coordinates so that the $`AdS_2`$ metric becomes $`ds^2=\mathrm{cos}^2\rho \left(dt^2+d\rho ^2\right),`$ $`\rho [\frac{\pi }{2},\frac{\pi }{2}]`$ and using the spectra of the Hamiltonians conjugate to the time variable $`t`$ calculated in  $`\omega _n^{(F)}(\mu )=n+|\mu |+\frac{1}{2},`$ $`\omega _n^{(B)}(\mu )=n+h(\mu ),h(\mu )=\frac{1}{2}(1+\sqrt{1+4m_B^2}),`$ one finds, summing over all the modes as in , the 1-loop vacuum energy of this effective 2-d field theory

$$E=\frac{1}{2}\underset{n=0}{\overset{\mathrm{}}{}}\left(3\left[\omega _n^{(B)}(1)\omega _n^{(F)}(1)\right]+5\left[\omega _n^{(B)}(1)\omega _n^{(F)}(1)\right]\right).$$
(13)
The properly defined vacuum energy should vanish in the $`AdS`$ case as it does in flat space  (even though divergences may not cancel out, unless there is a lot of supersymmetry ). The direct computation of the sum of the mode energies using $`\zeta `$-function regularization gives (using the standard relations $`\zeta (s,x)_{n=0}^{\mathrm{}}(n+x)^s,\zeta (1,x)=\frac{1}{2}\left(x^2x+\frac{1}{6}\right)`$)

$$E=\frac{1}{4}\left[3\times \left(2+\frac{1}{6}\right)+5\times \frac{1}{6}8\times \left(1\frac{1}{12}\right)\right]=0.$$
(14)
The ratio 3:5 of the numbers of the two multiplets is just what is needed for the cancellation. The fact that $`E`$, defined by $`\zeta `$-function regularization, vanishes should be a consequence of the extended $`𝒩=8`$ supersymmetry mentioned above. Indeed, while as was found for $`AdS_4`$ , the $`\zeta `$-function regularization may break supersymmetry and thus may lead to $`E0`$, this does not actually happen in the case of $`N5`$, $`d=4`$ gauged supergravities . The present $`d=2`$ case is thus analogous to those $`d=4`$ cases with large amounts of supersymmetry.

In curved (e.g., static conformally flat) space the logarithm of the partition function is, in general, different from the vacuum energy defined as a sum over eigen-modes because the time derivative part of the relevant elliptic operators is rescaled by $`g^{00}`$. The determinants of the two operators which differ by such a rescaling are related to each other by a conformal anomaly type equation. Taking this into account one finally concludes  that $`Z_{tot}=1`$. This result is a consequence of $`E=0`$ as well as of the cancellation of the sum of conformal anomalies for ten bosons with total mass terms 8, eight fermions with 4 times the standard 2-d fermion conformal anomaly and total mass 8, and the conformal gauge ghosts.

Similar results are found in the case of the string world sheet ending on a circular Wilson loop at the boundary, where the induced geometry is again $`AdS_2`$ .

## 4 One-loop correction to “bended” string
or rectangular Wilson loop

Analogous calculation of the one-loop correction can be done in the case of the “bended” string configuration with both ends being at the boundary of $`AdS_5`$ separated by distance $`L`$ in a spatial direction, and representing the rectangular $`(T,L)`$ Wilson loop in the boundary gauge theory .

The corresponding minimal surface was found in  and the classical value of the string action (equal to the minimal area as fermions vanish at the classical level) accounts for the leading large ‘t Hooft coupling $`\lambda =\frac{R^4}{\alpha ^2}`$ behavior ($`\frac{c_0\sqrt{\lambda }}{L}=\frac{c_0R^2}{\alpha ^{}L}`$) of the “quark – anti-quark” (W-boson) potential in the large $`N`$ $`𝒩=4`$ SYM theory. On general grounds, in this conformal field theory the Wilson loop should be $`<W(C)>`$exp$`[TV(L)]`$$`V(L)=\frac{f(\lambda )}{L}`$, where $`ff_0=d_1\lambda +d_2\lambda ^2+\mathrm{}`$ at weak coupling and AdS/CFT duality predicts  that $`ff_{\mathrm{}}=c_0\sqrt{\lambda }+c_1+O(\frac{1}{\sqrt{\lambda }})`$ at strong coupling. It is natural to expect that there is a smooth function $`f(\lambda )`$ which interpolates between the two expansions. The first correction $`\frac{c_1}{L}`$ to the strong-coupling expansion of the potential is determined by the one-loop ($`\alpha ^{}`$) correction to the string partition function (see also  for related discussions).

The use of the GS action of  as a starting point for this one-loop calculation was suggested in  (were the fermionic contribution was found in the gauge $`\theta ^1=\mathrm{\Gamma }_{0123}\theta ^2`$) and it was completed in .

Writing the $`AdS_5\times S^5`$ metric ($`R=1`$) as $`ds^2=y^2dx^ndx^n+\frac{dy^2}{y^2}+d\mathrm{\Omega }_5^2`$ where $`yx^4=w^1`$, the string configuration is described by

$$x^0=\tau (0,T),x^1=\sigma (\frac{L}{2},\frac{L}{2}),y=y(\sigma ),$$
(15)
where $`y(\sigma )`$ extremises the action  $`S=\frac{R^2}{2\pi \alpha ^{}}T𝑑\sigma \sqrt{y^2+y^4}`$, i.e. is determined by the second-order equation $`yy^{\prime \prime }=4y^2+2y^4`$ with the first integral being $`y^2=\frac{y^8}{y_0^4}y^4,`$ $`y_0=y_{min}=\frac{\kappa _0}{L},`$ $`\kappa _0\frac{(2\pi )^{3/2}}{[\mathrm{\Gamma }(\frac{1}{4})]^2}`$. The induced 2-d metric takes the form $`ds^2=y^2d\tau ^2+y^6d\sigma ^2`$, $`R^{(2)}=2(1+y^4)`$. The indiced geometry is asymptotic to $`AdS_2`$: if we change the coordinate $`\sigma `$ to $`y`$ the metric becomes

$$ds^2=y^2d\tau ^2+\frac{y^2}{y^4y_0^4}dy^2,y_0y<\mathrm{}.$$
(16)
The $`y_0=0`$ limit corresponds to the straight string configuration where the metric becomes that of Euclidean $`AdS_2`$ space (with $`0y<\mathrm{}`$).

In the flat space limit ($`R\mathrm{}`$) one finds that the quantum correction to the rectangular Wilson loop vanishes because of the cancellation of the bosonic and fermionic contributions due to effective 2-d supersymmetry present after gauge fixing. The 1-loop $`\frac{c_1}{L}`$ correction to the effective potential should not, however, vanish  in the present curved space case as there is no reason to expect that the action expanded near the solution $`y=y(\sigma )`$ should have an effective world-sheet supersymmetry.

The bosonic fluctuations near this bended string configuration were analysed in . As in the straight string case the one-loop bosonic partition function in the static gauge is again expressed in terms of massive and massless scalar determinants in induced 2-d geometry. In the static gauge

$$Z_B^{(stat.)}=\text{det}^{2/2}\left(^2+2\right)\text{det}^{1/2}\left(^2+R^{(2)}+4\right)\text{det}^{5/2}\left(^2\right),$$
(17)
while in the conformal gauge

$$Z_B^{(conf.)}=\frac{\text{det}^{1/2}\left(_{\mu \nu }^2\frac{1}{2}R^{(2)}g_{\mu \nu }\right)}{\text{det}^{1/2}\left(D_{ab}^2+X_{ab}\right)\text{det}^{5/2}\left(^2\right)},$$
(18)
where $`X_{ab}=2\delta _{ab}g^{\mu \nu }\eta _\mu ^a\eta _\nu ^b,`$ $`\eta _0^a=(y,0,0,0,0),`$ $`\eta _1^a=(0,y,0,0,y^1y^{})`$ (we set $`y_0=1`$). The two expressions are equivalent once supplemented with constant factors related to gauge fixing.

As discussed above, a systematic way to put the fermionic contribution into that of massive 2-d Dirac spinors in induced 2-d geometry is to perform a local Lorentz rotation. The quadratic part of the fermionic action depends on $`\rho _0=y\mathrm{\Gamma }_0,\rho _1=y\mathrm{\Gamma }_1+y^1y^{}\mathrm{\Gamma }_4,\widehat{D}_\mu =_\mu +\frac{1}{2}y\mathrm{\Gamma }_\mu \mathrm{\Gamma }_4.`$ Using the gauge $`\theta ^1=\theta ^2`$ (and Minkowski signature) one observes  that the combination of $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_4`$ which appears in $`\rho _1`$ can be interpreted as a (local, $`\sigma `$-dependent) rotation of $`\mathrm{\Gamma }_1`$ in the 1-4 plane $`S\mathrm{\Gamma }_1S^1=\mathrm{cos}\alpha \mathrm{\Gamma }_1+\mathrm{sin}\alpha \mathrm{\Gamma }_4=y^2\mathrm{\Gamma }_1+y^4y^{}\mathrm{\Gamma }_4=y^3\rho _1,`$ where $`S=\mathrm{exp}\left(\frac{\alpha }{2}\mathrm{\Gamma }_1\mathrm{\Gamma }_4\right),`$ $`\mathrm{cos}\alpha =y^2,`$ $`\alpha ^{}\frac{d\alpha }{d\sigma }=2y.`$ Making the field redefinition $`\theta \mathrm{\Psi }S^1\theta ,`$ one obtains

$$L_{2\mathrm{F}}=2i\sqrt{g}\left(\overline{\mathrm{\Psi }}\tau ^\mu \widehat{}_\mu \mathrm{\Psi }+i\overline{\mathrm{\Psi }}\tau _3\mathrm{\Psi }\right),\tau _3\frac{ϵ^{\mu \nu }}{2\sqrt{g}}\tau _\mu \tau _\nu =\mathrm{\Gamma }_0\mathrm{\Gamma }_1,(\tau _3)^2=1.$$
(19)
Choosing a representation for $`\mathrm{\Gamma }_a`$ such that $`\mathrm{\Gamma }_{0,1}`$ are $``$ 2-d Dirac matrices, i.e. $`\mathrm{\Gamma }_0=i\sigma _2\times I_8,\mathrm{\Gamma }_1=\sigma _1\times I_8,\tau _3=\mathrm{\Gamma }_0\mathrm{\Gamma }_1=\sigma _3\times I_8,`$ we end up with 8 species of 2-d Majorana fermions living on a curved 2-d surface with a $`\sigma _3`$ mass term. Assuming that fermions are normalized with $`\sqrt{g}`$, the square of the resulting fermionic operator is then $`\widehat{}^2+\frac{1}{4}R^{(2)}+1`$, where $`\widehat{}^2=\frac{1}{\sqrt{g}}\widehat{}^\mu (\sqrt{g}\widehat{}_\mu )`$ and $`\widehat{}`$ is the covariant 2-d spinor derivative of the induced metric. Equivalent results for the fermionic contribution were found in the “3-brane” gauge $`\theta ^1=i\mathrm{\Gamma }_4\theta ^2`$  and in a different $`\kappa `$-symmetry gauge in .

Combining the contributions of bosons in the static gauge and fermions in the $`\theta ^1=\theta ^2`$ gauge the total expression for the 1-loop partition function of a “bended” string in $`AdS_5\times S^5`$ is thus

$$Z=Z_0\frac{\text{det}^{8/2}(\widehat{}^2+\frac{1}{4}R^{(2)}+1)}{\text{det}^{2/2}\left(^2+2\right)\text{det}^{1/2}\left(^2+R^{(2)}+4\right)\text{det}^{5/2}\left(^2\right)},$$
(20)
where $`Z_0`$ is a gauge fixing (measure) factor which is the same as in the flat ($`R\mathrm{}`$) limit. In particular, this factor is crucial for cancellation  of the topological ($`\mathrm{log}ϵR^{(2)}`$) divergence present  in the ratio of the determinants. The key point is that the cancellation of this divergence in the one-loop approximation in curved target space is essentially the same as in the flat space, where to demonstrate that logarithmic Euler number divergences cancel (as they should to match the cancellation of conformal anomalies in $`D=10`$ superstring) one is to take into account various zero mode normalization factors in the string path integral measure .333One should note also that in performing a local target space rotation that transforms the quadratic GS fermion term into the 2-d fermion kinetic term the resulting Jacobian  depends on the 2-d metric and its contribution explains why the conformal anomaly of a GS fermion is 4 times bigger than that of a 2-d fermion. This is crucial for understanding how conformal anomalies cancel in $`D=10`$ GS string. To avoid questions about boundary terms (and details of topological infinity cancellation) one may normalize the partition function for each field by the partition function of an equivalent field in the straight string configuration, i.e. divide the partition function for the noncompact hyperbolic negative curvature 2-d space (16) by the partition function for the $`AdS_2`$ case.444 Since the topology and the near-boundary (large $`y`$) behavior of the two metrics is the same, this eliminates the problem of carefully tracking down all boundary terms in the expressions for the determinants and allows one to ignore the boundary contributions as well as the total derivative bulk terms (such as the logarithmically divergent terms proportional to $`d^2\sigma \sqrt{g}R^{(2)}`$). The ratio of the determinants for the metric (16) and for its $`y_0=0`$ limit will be finite and well-defined. This is the standard recipe of defining the determinants of Laplace operators on (e.g. 2-dimensional) non-compact spaces by using fiducial metrics of constant negative curvature which have the same topology and asymptotic behavior. Since the indiced 2-d geometry is static, the calculation of the partition function and thus of the one-loop coefficient $`c_1`$ in the potential $`V(L)`$ reduces to determining the spectra of one-dimensional Schrödinger operators, a tractable problem. A crude estimate for the resulting coefficient $`c_1`$ was discussed in .

## 5 Conclusions

To conclude, the Green-Schwarz action for the type IIB superstrings in $`AdS_5\times S^5`$ background provides a natural and well-defined starting point for developing $`\alpha ^{}`$ perturbation theory near “long” string configurations. The resulting partition function can be interpreted as that of an effective 2-d quantum field theory in curved induced 2-d geometry. The one-loop string correction is already sensitive to details of the string coupling to RR background and would be hard to reproduce in the standard NSR formalism where the RR vertex operator is defined near flat space. One may hope that progress towards understanding the spectrum of “short” closed strings in $`AdS_5\times S^5`$ and thus large $`N`$ superconformal YM theory for any value of $`\lambda `$ may be achieved using a light-cone gauge approach like the one initiated in .

Acknowledgements. I am grateful to R.R. Metsaev, N. Drukker and D. Gross for discussions and collaboration. This work was supported in part by the DOE grant DEFG02-91-ER-40690, the INTAS project 991590 and NATO grant PST.CLG 974965.


text:
# Stationary random fields with linear regressions Running head: Random fields with linear regressions

## 1 Introduction

The literature suggests that the distribution of a random sequence is often determined uniquely by the first two conditional moments. A classical result of P. Levy has been augmented by a number of results that replace assumptions on trajectories by more conditionings. Building on the work of Plucińska , Wesołowski  determined all stochastic processes with quadratic conditional variances under the assumptions that lead to processes with independent increments. Szabłowski ,  shows that first two conditional moments determine the marginal distribution in smooth cases. In  we show that the first two conditional moments characterize gaussian sequences. For more references and additional motivation the reader is referred to \[5, Chapter 8\].

In this paper we analyze stationary sequences under the assumptions which are invariant with respect to the ordering of the index set. Let $`(X_k)_k`$ be a square-integrable random field. We are interested in the fields with the first two conditional moments given by

$$E(X_k|\mathrm{},X_{k2},X_{k1},X_{k+1},X_{k+2},\mathrm{})=L(X_{k1},X_{k+1})$$
(1)
$$E(X_k^2|\mathrm{},X_{k2},X_{k1},X_{k+1},X_{k+2},\mathrm{})=Q(X_{k1},X_{k+1})$$
(2)
for all $`k`$ , where $`L(x,y)=a(x+y)+b`$ is a symmetric linear polynomial, and

$$Q(x,y)=A(x^2+y^2)+Bxy+C+D(x+y)$$
(3)
is a symmetric quadratic polynomial.

It turns out that conditions (1), (2) together with some technical assumptions imply that the eigenfunctions of the shifted conditional expectation are certain $`q`$-orthogonal polynomials, see (20). We explore this connection to identify the one-dimensional distributions of $`(X_k)`$ when the coefficients of the quadratic form (3) depend on one real parameter $`R`$ in addition to the correlation coefficient $`\rho `$.

We show that the following change occurs as the values of the parameter change.

* For a range of values between the two “critical values” of the parameter random field $`(X_k)`$ is bounded, its one-dimensional distribution is uniquely determined, and has a density.
* At the upper “critical value” of the parameter $`R`$ the distribution $`(X_0)`$ is normal.
* At the lower “critical value” of the parameter $`R`$ the distribution $`(X_0)`$ is two-valued.

The paper is organized as follows. Section 2 lists the assumptions used throughout the paper. Section 3 contains the statements of the main results. The next four sections contain lemmas and proofs. In section 8 we construct stationary Markov processes which satisfy (1) and (2). Section 9 collects concluding remarks.

## 2 Assumptions

The following general assumptions will be in place throughout the paper.

* We assume that $`(X_k)`$ is $`L_2`$-stationary and that all one dimensional distributions are equal,

$$X_k\stackrel{~}{=}X_0\text{ for all }k$$
(4)
* Since (1) and (2) are not affected by non-degenerate linear transformations of $`X_k`$, without loss of generality we assume $`E(X_k)=0`$ and $`E(X_k^2)=1`$. All specific assumptions about the coefficients of $`Q(x,y)`$ refer to this case.
* We put the following restrictions on the values of the first two correlation coefficients $`\rho =corr(X_0,X_1)`$ and $`r_2=corr(X_0,X_2)`$.

$$\rho 0$$
(5)
$$r_2+12|\rho |>0$$
(6)

###### Remark 2.1

Note that (5) excludes i. i. d. sequences, and (6) implies $`|\rho |<1`$ thus excluding sequences formed by repeating a single random variable. Clearly, both i. i. d. sequences and constant sequences with second moments satisfy (1) and (2) but can have arbitrary distribution. It is somewhat surprising that for other values of $`\rho `$ the distribution of $`X_0`$ is determined uniquely.

## 3 Theorems

Under mild assumptions condition (1) determines the covariances of $`(X_k)`$ up to a multiplicative factor.

###### Theorem 3.1

Let $`(X_k)`$ be an $`L_2`$-stationary random field such that (1) holds true for all $`k`$, and conditions (5) and (6) are satisfied. Let $`r_k=corr(X_0,X_k)`$ denote the correlation coefficients.

Then

* $`r_k=\rho ^{|k|}`$ for all $`k`$.
* One-sided regressions are also linear,

$$E(X_k|\mathrm{},X_1,X_0)=r_kX_0,k=1,2\mathrm{}$$
(7)

By symmetry, Theorem 3.1 implies the following.

###### Corollary 3.1

Under the assumptions of Theorem 3.1, $`E(X_k|X_0)=\rho ^{|k|}X_0`$ for all $`k`$.

We will be interested in how the properties of one dimensional distribution $`(X_0)`$ change with the values of the coefficients of the quadratic form $`Q`$.

###### Proposition 3.1

Let $`(X_k)`$ be an $`L_2`$-stationary standardized random field with conditional moments given by (1) and (2), and such that conditions (5-6) are satisfied. Suppose that (3) holds with $`A<1/(1+\rho ^2)`$. If $`A(\rho ^2+\frac{1}{\rho ^2})+B<1`$ then $`X_0`$ is bounded. If $`A(\rho ^2+\frac{1}{\rho ^2})+B=1`$ and $`D0`$ then either $`X_0`$ is bounded from below ($`D>0`$) or from above ($`D<0`$).

Proposition 3.1 suggests that the case when

$$A(\rho ^2+\frac{1}{\rho ^2})+B=1$$
(8)
and $`D=0`$ is of interest. This case includes the Gaussian fields, where $`A=\rho ^2/(1+\rho ^2)^2`$ and $`B=2\rho ^2/(1+\rho ^2)^2`$; in fact, we show in  that for a Markov chain which satisfies (1) and (2), the coefficients of (3) either satisfy constraint (8), or $`2A+B\rho ^2=1`$, and that $`D=0`$. Since by taking the expected value of (2) we get a trivial relation

$$C=12AB\rho ^2$$
(9)
(here we used $`r_2=\rho ^2`$ by Theorem 3.1) this leaves just one free parameter besides $`\rho `$ on the right hand side of (3). As this free parameter we will use a scaled version of $`B`$ defined by

$$R=B(\rho +\frac{1}{\rho })^2$$
(10)
###### Theorem 3.2

Let $`(X_k)`$ be an $`L_2`$-stationary standardized random field with conditional moments given by (1) and (2), and such that conditions (4-6) are satisfied. Furthermore, assume that the coefficients on the right hand side of (3) are such that $`D=0`$ and (8) holds true. Then the following three statements hold.

* For $`0R<2`$ the one-dimensional distribution $`(X_0)`$ has the uniquely determined symmetric distribution supported on a finite interval. When $`R=0`$, $`X_0=\pm c`$ has two values.
* For $`R=2`$ the one-dimensional distribution $`(X_0)`$ is normal.
* For $`R>2`$ the one-dimensional distributions $`(X_0)`$ are not determined by moments.

The infinite number of variables in the conditional expectations in (1) and (2) cannot be easily reduced.

###### Theorem 3.3

For every value of $`1<\rho <1`$ there is a non-degenerate $`L_2`$-stationary random field which satisfies condition (4), has linear regressions

$$E(X_k|X_{k1},X_{k+1})=\alpha (X_{k1}+X_{k+1}),$$
(11)
has quadratic second moments

$$E(X_k^2|X_{k1},X_{k+1})=\alpha ^2(X_{k1}+X_{k+1})^2+C,$$
(12)
and

* the conclusion of Theorem 3.1(i) fails;
* the conclusion of Theorem 3.2 fails.

In Section 8, for every value of parameters $`0<R2`$ and $`0<|\rho |<1`$ we construct strictly stationary Markov chains which satisfy the assumptions of Theorem 3.2.

## 4 Proof of Theorem 3.1

We use the notation $`E(|\mathrm{},X_n)`$ to denote the conditional expectation with respect to the sigma field generated by $`\{X_k:kn\}`$. Occasionally, we will also write $`E^{}(X):=E(X|)`$.

Let $`r_k=E(X_0X_k)`$ be the correlation coefficients, $`k=0,1,\mathrm{}`$, and recall that $`\rho =r_10`$ by assumption.

Theorem 3.1 follows from the following two observations.

###### Lemma 4.1

Under the assumptions of Theorem 3.1

* $`r_k0`$ as $`k\mathrm{}`$.
* $`r_k^2<1`$ for all $`k>2`$.

Proof. Multiplying (1) by $`X_0`$ and taking the expected value of both sides we get $`r_k=a(r_{k1}+r_{k+1})`$. Since $`\rho 0`$, (6) implies that $`1+r_2>0`$. Therefore $`a=\rho /(1+r_2)`$ and the correlation coefficients $`r_k`$ satisfy the recurrence

$$(1+r_2)r_k=\rho (r_{k1}+r_{k+1}),k=1,2,\mathrm{}$$

For fixed $`\rho ,r_2`$, this is a linear recurrence for $`(r_k)_{k>2}`$. Inequality (6) implies that the characteristic equation of the recurrence has two distinct real roots and their product is $`1`$. There is therefore only one root, $`c`$, in the interval $`(1,1)`$ and $`r_k=bc^k`$. This shows that $`r_k=r_2c^{k2},k=2,3,\mathrm{}`$. In particular, $`r_k0`$, and $`r_k^2<1`$ for $`k>2`$.

###### Lemma 4.2

If $`(X_k)`$ satisfies the assumptions of Theorem 3.1, then

$$E(X_{n+1}|\mathrm{},X_n)=\rho X_n$$
(13)
Proof. We first show that for all $`n,k2`$ there are coefficients $`a(k)0`$ and $`b(k)`$ such that

$$E(X_{n+1}|\mathrm{},X_{n1},X_n,X_{n+k},X_{n+k+1},\mathrm{})=a(k)X_n+b(k)X_{n+k}$$
(14)
and

$$E(X_{n+k1}|\mathrm{},X_{n1},X_n,X_{n+k},X_{n+k+1},\mathrm{})=b(k)X_n+a(k)X_{n+k}.$$
(15)
We prove this by induction with respect to $`k2`$.

For $`k=2`$, this follows from (1) with $`a(2)=\frac{\rho }{1+r_2}0`$, see (5).

For a given $`k2`$, suppose that $`a(k)0`$ and both (14) and (15) hold true for all $`n`$. We will prove that the same statement holds true for $`k+1`$.

Conditioning on additional variable $`X_{n+k}`$ and using the induction assumption we get

$$E(X_{n+1}|\mathrm{},X_{n1},X_n,X_{n+k+1},X_{n+k+2},\mathrm{})=$$

$$E^{\mathrm{},X_{n1},X_n,X_{n+k+1},\mathrm{}}\left(E^{\mathrm{}X_{n1},X_n,X_{n+k},X_{n+k+1},\mathrm{}}(X_{n+1})\right)$$

$$a(k)X_n+b(k)E(X_{n+k}|\mathrm{},X_{n1},X_n,X_{n+k+1},\mathrm{})$$

Now adding $`X_{n+1}`$ to the condition in $`E(X_{n+k}|\mathrm{},X_{n1},X_n,X_{n+k+1},\mathrm{})`$ and using the induction assumption we get

$$E(X_{n+k}|\mathrm{},X_{n1},X_n,X_{n+k+1},X_{n+k+2},\mathrm{})=$$

$$E^{\mathrm{},X_{n1},X_n,X_{n+k+1},X_{n+k+2},\mathrm{}}\left(E^{\mathrm{},X_n,X_{n+1},X_{n+k+1},X_{n+k+2},\mathrm{}}(X_{n+k})\right)=$$

$$b(k)E(X_{n+1}|\mathrm{},X_{n1},X_n,X_{n+k+1},X_{n+k+2},\mathrm{})+a(k)X_{n+k+1}$$

Combining these two expressions we get the linear equation

$$M=a(k)X_n+b^2(k)M+a(k)b(k)X_{n+k+1}$$

for unknown random variable $`M=E(X_{n+1}|\mathrm{},X_{n1},X_n,X_{n+k+1},X_{n+k+2},\mathrm{})`$. Notice that since $`a(k)0`$, if $`b^2(k)=1`$ then $`X_n,X_{n+k+1}`$ are linearly related. Therefore, $`r_{k+1}^2=1`$, contradicting Lemma 4.1(ii). (Case $`k=2`$ requires separate verification: $`b^2(2)=\frac{\rho ^2}{(1+r_2)^2}<1`$ by (6)). Thus $`b^2(k)1`$ and $`M`$ is determined uniquely as the linear function of $`X_n,X_{n+k+1}`$. This proves (14). By symmetry (or by similar reasoning) (15) holds true. The equation gives $`a(k+1)=\frac{a(k)}{1b^2(k)}`$ which shows that $`a(k+1)0`$. This completes the induction proof of (14).

Since $`|r_k|<1`$ for $`k>2`$, the coefficients $`a(k),b(k)`$ in a linear regression are determined uniquely. A calculation gives $`b(k)=\frac{r_{k1}r_1r_k}{1r_k^2}`$ Using Lemma 4.1(i) we have $`b(k)0`$ as $`k\mathrm{}`$. Passing to the limit as $`k\mathrm{}`$ in (14) we get (13).

Proof of Theorem 3.1. Applying $`k`$-times formula (13), we get $`E(X_k|\mathrm{},X_0)=\rho ^kX_0`$. This implies $`E(X_0X_k)=\rho ^k`$.

## 5 Proof of Proposition 3.1

A simple calculation using (8) and (10) gives the following.

###### Lemma 5.1

The following two conditions are equivalent:

* $`A<1/(1+\rho ^2)`$;
* $`R>11/\rho ^4`$.

The second part of the next lemma will serve as the first step in the induction proof of Lemma 6.3.

###### Lemma 5.2

If $`(X_k)`$ satisfies the assumptions of Theorem 3.1 and (2) holds true, then

$$(1A(1+\rho ^2))E(X_1^2|\mathrm{},X_0)=(A(1\rho ^2)+B\rho ^2)X_0^2+C+D(1+\rho ^2)X_0.$$
(16)
If in addition (8) holds true and $`R0`$ then

$$E(X_1^2|X_0)=\rho ^2X_0^2+1\rho ^2+\gamma X_0,$$
(17)
where $`\gamma =\frac{D}{1(1+\rho ^2)A}`$.

Proof. Since $`E(X_1X_2|\mathrm{},X_0)=E^{\mathrm{},X_0}\left(X_1E^{\mathrm{},X_1}(X_2)\right)`$, from Lemma 4.2 we get

$$E(X_1X_2|\mathrm{},X_0)=\rho E(X_1^2|\mathrm{},X_0).$$
(18)
We now give another expression for the left hand side of (18). Theorem 3.1 implies that $`L(x,y)=\frac{\rho }{1+\rho ^2}(x+y)`$. Since $`E(X_1X_2|\mathrm{},X_0)=E(X_2E(X_1|\mathrm{},X_0,X_2,\mathrm{})|\mathrm{},X_0)`$ from (1) we get $`E(X_1X_2|\mathrm{},X_0)=\frac{\rho }{1+\rho ^2}E(X_2(X_2+X_0)|\mathrm{},X_0)`$. By Lemma 4.2 this implies $`E(X_1X_2|\mathrm{},X_0)=\frac{\rho ^3}{1+\rho ^2}X_0^2+\frac{\rho }{1+\rho ^2}E(X_2^2|\mathrm{},X_0)`$. Since $`\rho 0`$, combining the latter with (18) we have

$$E(X_2^2|\mathrm{},X_0)=(1+\rho ^2)E(X_1^2|\mathrm{},X_0)\rho ^2X_0^2.$$
(19)
We now substitute expression (19) into (2) as follows. Taking the conditional expectation $`E(|\mathrm{},X_0)`$ of both sides of (2) with $`k=1`$ we get

$$E(X_1^2|\mathrm{},X_0)=AX_0^2+AE(X_2^2|\mathrm{},X_0)+BX_0^2\rho ^2+C+D(1+\rho ^2)X_0.$$

Replacing $`E(X_2^2|\mathrm{},X_0)`$ in the above expression by the right hand side of (19) we get (16).

If (8) holds true then identity (17) follows by a simple calculation, since by Lemma 5.1 and the assumption $`R0`$ we have $`A<1/(1+\rho ^2)`$.

By symmetry, we can switch the roles of $`X_0,X_1`$ in (17). Therefore the assumptions of \[4, Corollary 3\] are satisfied. This implies integrability, see also \[16, Theorem 2\] for another proof.

###### Corollary 5.1

Under the assumptions of Theorem 3.1 $`E(|X_k|^p)<\mathrm{}`$ for all $`p>1`$.

### 5.1 Proof of Proposition 3.1

By Theorem 3.1(ii) and (16)

$$\rho ^2X_0^2=(E^{X_0}(X_1))^2E^{X_0}(X_1^2)=\frac{A(1\rho ^2)+B\rho ^2}{1(1+\rho ^2)A}X_0^2+\frac{D}{1(1+\rho ^2)A}X_0+\frac{C}{1(1+\rho ^2)A}.$$

To prove the first part, notice that if $`X_0`$ is unbounded then $`\rho ^2\frac{A(1\rho ^2)+B\rho ^2}{1(1+\rho ^2)A}`$. The latter implies that $`A(\rho ^2+1/\rho ^2)+B1`$. Indeed, write $`A(\rho ^2+1/\rho ^2)+B=1+\mathrm{\Delta }`$. Then $`\frac{A(1\rho ^2)+B\rho ^2}{1(1+\rho ^2)A}=\rho ^2+\frac{\rho ^2\mathrm{\Delta }}{1(1+\rho ^2)A}`$, so $`\mathrm{\Delta }0`$.

Furthermore, if (8) holds true then an elementary calculation gives $`\frac{A(1\rho ^2)+B\rho ^2}{1(1+\rho ^2)A}=\rho ^2`$. Thus condition $`D>0`$ implies that $`X_0`$ must be bounded from below. Similarly, if $`D<0`$ then $`X_0`$ must be bounded from above.

## 6 Proof of Theorem 3.2

### 6.1 Orthogonal polynomials

Define polynomials $`Q_n`$ by the recurrence

$$xQ_n(x)=Q_{n+1}(x)+(1+q+\mathrm{}+q^{n1})Q_{n1}(x)$$
(20)
with $`Q_1(x)=0,Q_0(x)=1`$. (Then $`Q_1(x)=x`$ and $`Q_2(x)=x^21`$). Clearly, $`Q_n`$ are the Hermite polynomials when $`q=1`$.

Polynomials $`Q_n`$ after a change of variable transform into the continuous $`q`$-Hermite polynomials, see  and . Thus we can deduce information about the measure that makes $`Q_n`$ orthogonal; this measure is unique (with explicit representation for the density) and has bounded support when $`1<q<1`$, normal when $`q=1`$, and is non-unique when $`q>1`$. Later on we shall see that $`Q_n`$ are orthogonal with respect to the distributions of $`(X_0)`$.

The next lemma collects these results together. A sketch of the proof is enclosed for completeness.

###### Lemma 6.1

Let $`Q_n`$ be defined by (20).

* If $`q<1`$ then there is no probability measure with respect to which $`Q_n`$ are orthogonal.
* If $`1q<1`$ then $`Q_n`$ are orthogonal with respect to a unique probability measure. This measure is symmetric on $``$ and has bounded support.
* If $`q=1`$ then $`Q_n`$ are orthogonal with respect to a unique probability measure which is normal.
* If $`q>1`$ then $`Q_n`$ are orthogonal with respect to an infinite number of probability measures.

Proof. (i) Follows from the fact that the coefficients in (20) must be non-negative.

(ii) This is an immediate consequence of Carleman’s criterion, see \[13, page 59\].

It is easy to see that $`q=1`$ corresponds to the (unique) symmetric two-point distribution $`X_0=\pm 1`$. This is a degenerate case, often excluded from the general theory of orthogonal polynomials because $`Q_n(X_0)=0`$ for all $`n>1`$.

(iii) This is the consequence of the classical result that normal distribution is determined uniquely by its moments.

(iv) This follows from Berezanskii’s result as reported in addendum 5 (page 26) of . Notice that the latter deals with normalized rather than monic polynomials, see (30) below. Since $`E(Q_n^2(X_0))=_{k=1}^n\frac{q^k1}{q1}`$, Berezanskii’s theorem is used with $`b_n=\sqrt{\frac{q^{n+1}1}{q1}}`$ and the condition $`b_{n+1}b_{n1}b_n^2`$ holds true because $`x\mathrm{log}(q^x1)`$ is concave when $`q>1`$.

Several explicit weight functions for $`q`$-Hermite polynomials with $`q>1`$ are given in .

The following lemma shows that we can find orthogonal polynomials by finding the eigenfunctions of conditional expectations, and is based on well known motives.

###### Lemma 6.2

Let $`f,g`$ be two functions such that $`f(X_1)`$ and $`g(X_1)`$ are square-integrable. If $`(X_0)=(X_1)`$ and $`\alpha \beta `$ are real numbers such that $`E(f(X_1)|X_0)=\alpha f(X_0)`$ and $`E(g(X_0)|X_1)=\beta g(X_1)`$ then $`E(f(X_0)g(X_0))=0`$.

Proof. This follows from $`E(f(X_1)g(X_0))=\alpha E(f(X_0)g(X_0))=\beta E(f(X_1)g(X_1))`$.

### 6.2 Proof of Theorem 3.2

Let

$$q=\frac{\rho ^4+R1}{1+\rho ^4(R1)}.$$
(21)
It is easy to see that the range $`0R2`$ corresponds to $`1q1`$ and that $`R=2`$ when $`q=1`$, and $`R=0`$ when $`q=1`$.

###### Lemma 6.3

If $`q`$ is defined by (21) then polynomials $`Q_n(x)`$ defined by (20) satisfy

$$E^{X_0}Q_n(X_1)=\rho ^nQ_n(X_0).$$
(22)
Proof. The proof is by mathematical induction with respect to $`n`$. The result is trivially true for $`n=1`$ (with $`Q_1=0`$) and for $`n=0`$. By Theorem 3.1, (22) holds true for $`n=1`$. Since we assume that $`D=0`$, formula (17) implies that (22) holds true for $`n=2`$.

Suppose (22) holds true for $`n,n1,n2`$, and $`n3`$. We will show that (22) holds for $`n+1`$.

The proof repeatedly uses conditional expectations $`()=E^{\mathrm{},X_1,X_0}()`$.

Consider $`(X_1Q_n(X_2))`$. We will first obtain two equations (23), (24) which allow us to determine conditional moments $`(X_1Q_n(X_2))`$ and $`(X_1Q_n(X_1))`$.

The first equation is obtained as follows. By induction assumption, $`(X_1Q_n(X_2))=(X_1E^{\mathrm{},X_0,X_1}Q_n(X_2))=\rho ^n(X_1Q_n(X_1))`$. On the other hand, equation (1) gives

$$(X_1Q_n(X_2))=\frac{\rho }{1+\rho ^2}((X_0Q_n(X_2))+(X_2Q_n(X_2))).$$

Therefore

$$(X_2Q_n(X_2))=(1+\rho ^2)\rho ^{n1}(X_1Q_n(X_1))\rho ^{2n}X_0Q_n(X_0).$$
(23)
To obtain the second equation, consider $`(X_1^2Q_{n1}(X_2))`$. By induction assumption, $`(X_1^2Q_{n1}(X_2))=\rho ^{n1}(X_1^2Q_{n1}(X_1))`$. On the other hand, equations (2) and (3) with $`D=0`$ give

$$(X_1^2Q_{n1}(X_2))=$$

$$A(X_2^2Q_{n1}(X_2))+A(X_0^2Q_{n1}(X_2))+$$

$$B(X_0X_2Q_{n1}(X_2))+C(Q_{n1}(X_2)).$$

(We will use relations (8) and (9) later on.)

Thus, using the induction assumption, we obtain the second equation

$`\rho ^{n1}(X_1^2Q_{n1}(X_1))=`$ (24)
$`A(X_2^2Q_{n1}(X_2))+A\rho ^{2n2}X_0^2Q_{n1}(X_0)+`$
$`BX_0(X_2Q_{n1}(X_2))+C\rho ^{2n2}Q_{n1}(X_0).`$
Polynomials $`Q_n`$ satisfy second order recurrence of the form $`xQ_n(x)=Q_{n+1}(x)+\beta _{n+1}Q_{n1}(x)`$. We use it to rewrite (24) as follows.

$`\rho ^{n1}(X_1Q_n(X_1))+\rho ^{n1}\beta _n(X_1Q_{n2}(X_1))=`$ (25)
$`A(X_2Q_n(X_2))+A\beta _n(X_2Q_{n2}(X_2))+A\rho ^{2n2}X_0Q_n(X_0)+A\rho ^{2n2}\beta _nX_0Q_{n2}(X_0)+`$
$`BX_0(Q_n(X_2))+BX_0\beta _n(Q_{n2}(X_2))+C\rho ^{2n2}Q_{n1}(X_0).`$
Now use (23) to eliminate $`(X_2Q_n(X_2))`$ from this equation. Since $`C=12A\rho ^2B`$ this gives

$`(1A(1+\rho ^2))(X_1Q_n(X_1))=`$ (26)
$`\beta _n(A(1+\rho ^2)\rho ^21)(X_1Q_{n2}(X_1))+`$
$`\beta _n\rho ^{n3}(A(\rho ^21)+B)X_0Q_{n2}(X_0)+`$
$`\rho ^{n1}(A(1\rho ^2)+B\rho ^2)X_0Q_n(X_0)+`$
$`(12A\rho ^2B)\rho ^{n1}Q_{n1}(X_0).`$
By (8), identity $`\frac{A(\rho ^2+1)\rho ^2}{\rho ^2(1A(1+\rho ^2))}=\rho ^2q`$ (recall that $`q`$ is given by (21)) and elementary calculation we now get

$`(X_1Q_n(X_1))=`$ (27)
$`\beta _nq\rho ^2\left(\rho ^{n3}X_0Q_{n2}(X_0)(X_1Q_{n2}(X_1))\right)+`$
$`\rho ^{n+1}X_0Q_n(X_0)+(1\rho ^2)\rho ^{n1}Q_{n1}(X_0)`$
Since $`xQ_n(x)=Q_{n+1}(x)+\beta _{n+1}Q_{n1}(x)`$, and by induction assumption $`(Q_{n3}(X_1))=\rho ^{n3}Q_{n3}(X_0)`$, this implies

$`(Q_{n+1}(X_1))=\rho ^{n+1}Q_{n+1}(X_0)+`$
$`(1\rho ^2)\rho ^{n1}\left(q\beta _n\beta _{n+1}+1\right)Q_{n1}(X_0).`$
Since (20) means that $`\beta _{n+1}=q\beta _n+1`$, this ends the proof.

Proof of Theorem 3.2. All assumptions are symmetric with respect to time-reversal. Therefore, formula (22) implies $`E(Q_n(X_0)|X_1)=\rho ^nQ_n(X_1)`$. Since $`0<|\rho |<1`$, Lemma 6.2 (used with $`f=Q_m,g=Q_n,nm`$) proves that $`Q_n`$ are orthogonal with respect to $`(X_0)`$. Therefore, Lemma 6.1 identifies uniquely the distribution $`(X_0)`$ for both $`R=2`$ and $`R<2`$ cases.

The distribution is not determined uniquely by the moments when $`q>1`$ which corresponds to $`R>2`$.

Finally, when the distribution is determined uniquely, the odd-order moments $`EX_0^{2n+1}=0`$ by (20). Therefore the distribution of $`X_0`$ is symmetric.

## 7 Proof of Theorem 3.3

We begin with the following simple lemma.

###### Lemma 7.1

Let $`X_0,Y_0`$ be integrable random variables. Suppose that $`,𝒢`$ are $`\sigma `$-fields such that $`\sigma (X_0,)`$ and $`\sigma (Y_0,𝒢)`$ are independent, and there is $`\rho `$ such that $`E(X_0|)=\rho X_1`$, $`E(Y_0|𝒢)=\rho Y_1`$, $`E(X_0^2|)=\rho ^2X_1^2+1\rho ^2`$ , and $`E(Y_0^2|𝒢)=\rho ^2Y_1^2+1\rho ^2`$.

Let $`Z_k=(aX_k+bY_k)`$, where $`a^2+b^2=1`$ and denote by $`𝒢`$ the $`\sigma `$-field generated by $`𝒢`$ . Then

$$E(Z_0|𝒢)=\rho Z_1,$$
(28)
$$E(Z_0^2|𝒢)=\rho ^2Z_1^2+1\rho ^2.$$
(29)
Proof. Clearly $`E(Z_0|𝒢)=aE(X_0|𝒢)+bE(Y_0|𝒢)`$ proving (28). Similarly, $`E(Z_0^2|𝒢)=a^2E(X_0^2|𝒢)+b^2E(Y_0^2|𝒢)+2abE(Y_0X_0|𝒢)`$. Now $`E(Y_0X_0|𝒢)=E(Y_0E(X_0|𝒢,Y_0)|𝒢)=\rho E(Y_0X_1|𝒢)`$. Since $`X_1`$ is $`𝒢`$-measurable, we get $`E(Y_0X_0|𝒢)=\rho ^2X_1Y_1`$, which proves (29).

### 7.1 Proof of Theorem 3.3

Fix $`1<\rho <1`$. First define a periodic stationary sequence $`(\xi _k)`$ such that $`\xi _{k+2}=\xi _k`$ and with the correlation coefficient $`\rho `$. To this end define the joint distribution of $`\xi _1,\xi _2`$ by

$$\mathrm{Pr}(\xi _1=1,\xi _2=1)=\mathrm{Pr}(\xi _1=1,\xi _2=1)=\frac{1+\rho }{4},$$

$$\mathrm{Pr}(\xi _1=1,\xi _2=1)=\mathrm{Pr}(\xi _1=1,\xi _2=1)=\frac{1\rho }{4}.$$

Let $`(\gamma _k)`$ be a centered Markov gaussian sequence with correlations $`E(\gamma _0\gamma _k)=r^k`$, where $`r=\frac{1\sqrt{1\rho ^2}}{\rho }`$, and independent of $`(\xi _k)`$.

Let $`X_k=a\xi _k+b\gamma _k`$, where $`a^2+b^2=1`$. Then $`E(X_0X_1)=a^2\rho +b^2r`$. By selecting $`a`$ close enough to $`1`$, and by varying $`\rho `$ we can thus have correlations $`corr(X_0,X_1)`$ fill out the entire interval $`(1,1)`$.

Using Lemma 7.1 we verify that $`(X_k)`$ satisfies (11) and (12) with $`\alpha =\rho /2,C=1\rho ^2`$. Indeed, $`E(\xi _1|\xi _0,\xi _2)=\rho \xi _0=\rho (\xi _0+\xi _2)/2`$ and $`E(\xi _1^2|\xi _0,\xi _2)=1=\rho ^2(\xi _0+\xi _2)^2/4+1\rho ^2`$. Similarly, the gaussian sequence satisfies $`E(\gamma _1|\gamma _0,\gamma _2)=\frac{r}{1+r^2}(\gamma _0+\gamma _2)=\rho (\gamma _0+\gamma _2)/2`$ and $`E(\gamma _1^2|\gamma _0,\gamma _2)=\rho ^2(\gamma _0+\gamma _2)^2/4+1\rho ^2`$

When $`a0`$ the correlation coefficients $`E(X_0X_k)`$ do not converge to $`0`$. Thus the conclusion of Theorem 3.1(ii) fails.

Notice that $`Q=L^2+const`$ which corresponds to the normal case ($`R=0`$); the distribution of $`X_0`$ is not compactly-supported but for $`a0`$ it is not normal. Thus the conclusion of Theorem 3.2 fails.

## 8 Construction of Markovian fields

In this section we construct Markov fields which satisfy (1) and (2), proving the following.

###### Proposition 8.1

For all $`0<\rho ^2<1,0R2`$, there are stationary Markov chains that satisfy the assumptions of Theorem 3.2.

Since $`R=0`$ corresponds to the elementary two-valued Markov chain which was explicitly analyzed in , and $`R=2`$ corresponds to stationary Gaussian Markov processes, we only consider the case $`0<R<2`$.

Through the remainder of this section we fix $`\rho ,R`$. Let $`Q_n`$ be given by the recurrence

$$xQ_n(x)=b_{n+1}Q_{n+1}(x)+b_nQ_{n1}(x)$$
(30)
with initial polynomials $`Q_0(x)=1,Q_1(x)=x`$, where $`b_n=\sqrt{1+q+\mathrm{}+q^{n1}},b_0=0,b_1=1`$, and $`q`$ is given by (21). Since $`q>1`$, $`b_n>0`$ for all $`n1`$. These are the normalized orthogonal polynomials from the proof of Theorem 3.2.

Let $`\mu `$ be a probability measure which makes $`Q_n`$ orthogonal.

###### Lemma 8.1

If $`1<q1`$ and $`|\rho |<1`$, then $`_{n=0}^{\mathrm{}}\rho ^nQ_n(y)Q_n(x)`$ converges in $`L_2(\mu (dx)\mu (dy))`$ to a non-negative function.

Proof. The case $`q=1`$ is classical and the sum of the series is $`\frac{1}{\sqrt{1\rho ^2}}\mathrm{exp}(\rho \frac{\rho (x^2+y^2)2xy}{2(1\rho ^2)})`$.

For $`1<q<1`$ an explicit product representation for the series can be deduced from the facts collected in , see also \[8, (3.9)\]. Namely, in the notation of , $`Q_n(x)=\frac{e^{in\theta }H_n(e^{2in\theta })}{(1q)^{n/2}b_1\mathrm{}b_n}`$ where $`x=2\mathrm{cos}(\theta )/\sqrt{1q}`$. Therefore, using \[8, (3.9)\] with $`x=2\mathrm{cos}(\theta _x)/\sqrt{1q}`$, $`y=2\mathrm{cos}(\theta _y)/\sqrt{1q}`$ (see also \[9, (2.10), (5.2), and (5.7)\]) we have

$$\underset{n=0}{\overset{\mathrm{}}{}}\rho ^nQ_n(y)Q_n(x)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1\rho ^2q^k)}{(1+\rho ^2q^{2k}2\rho q^k\mathrm{cos}(\theta _x+\theta _y))(1+\rho ^2q^{2k}2\rho q^k\mathrm{cos}(\theta _x\theta _y))}.$$

Since the last expression is a product of positive factors, this ends the proof when $`|q|<1`$.

Let $`(X_n)`$ be a Markov chain with initial distribution $`\mu `$ and the transition probability

$$P_x(dy)=\underset{n=0}{\overset{\mathrm{}}{}}\rho ^nQ_n(y)Q_n(x)\mu (dy).$$

###### Lemma 8.2

$`(X_n)`$ is stationary, and satisfies condition (1).

Proof. Applying Fubini’s theorem to the function which by Lemma 8.1 is non-negative we get $`P_x(A)\mu (dx)=_A_R_{n=0}^{\mathrm{}}\rho ^nQ_n(y)Q_n(x)\mu (dx)\mu (dy)`$. Since $`Q_n(x)\mu (dx)=0`$ for all $`n>0`$ and $`Q_0=1`$, we get $`P_x(A)\mu (dx)=\mu (A)`$. This proves stationarity.

To prove (1), by the Markov property it suffices to show that $`E(X_1|X_0,X_2)=\frac{\rho }{1+\rho ^2}(X_0+X_2)`$. Let $`\varphi (X_0,X_2)0`$ be an arbitrary bounded measurable function. We will verify that

$$E(X_1\varphi (X_0,X_2))=\frac{\rho }{1+\rho ^2}E((X_0+X_2)\varphi (X_0,X_2)).$$
(31)
Since $`\varphi `$ is square integrable, and $`Q_{n,m}(x,y)=Q_n(x)Q_m(y)`$ are orthogonal in $`L_2(\mu (dx)\mu (dy))`$ we can write $`\varphi (X_0,X_2)=\varphi _0(X_0,X_2)+_{i,j=0}^{\mathrm{}}\varphi _{i,j}Q_i(X_0)Q_j(X_2)`$, where $`\varphi _0`$ is orthogonal to all polynomials in variables $`X_0,X_2`$. In the case $`R<2`$ we have $`\varphi _0=0`$ because $`(Q_n)`$ are an orthogonal basis of $`L_2(\mu )`$.

Notice that the joint distribution of $`X_0,X_1,X_2`$ is given by

$$\nu (dx,dy,dz)=\underset{n,k=0}{\overset{\mathrm{}}{}}\rho ^{n+k}Q_n(x)Q_n(y)Q_k(y)Q_k(z)\mu (dx)\mu (dy)\mu (dz).$$
(32)
From (30) and orthogonality we get

$$yQ_n(y)Q_{n1}(y)\mu (dy)=b_n$$

Moreover, $`yQ_n(y)Q_k(y)\mu (dy)=0`$ for $`kn\pm 1`$.

Therefore

$$y\varphi (x,z)\nu (dx,dy,dz)=\underset{n,k}{}\rho ^{n+k}\varphi _{n,k}yQ_n(y)Q_k(y)\mu (dy)$$

$$=\rho \rho ^{2n}b_{n+1}(\varphi _{n,n+1}+\varphi _{n+1,n}).$$

Similar calculation gives

$$x\varphi (x,z)\nu (dx,dy,dz)=\rho ^2\underset{n=0}{\overset{\mathrm{}}{}}\rho ^{2n}b_{n+1}\varphi _{n,n+1}+\underset{n=0}{\overset{\mathrm{}}{}}\rho ^{2n}b_{n+1}\varphi _{n+1,n}.$$

Since by symmetry a similar formula holds true for the integral of $`z`$ instead of $`x`$, we get

$$(x+z)\varphi (x,z)\nu (dx,dy,dz)=(1+\rho ^2)\underset{n=0}{\overset{\mathrm{}}{}}\rho ^{2n}b_{n+1}(\varphi _{n,n+1}+\varphi _{n+1,n}).$$

Comparing the coefficients in the expansions we verify that

$$y\varphi (x,z)\nu (dx,dy,dz)=\frac{\rho }{1+\rho ^2}(x+z)\varphi (x,z)\nu (dx,dy,dz),,$$

proving (31).

###### Lemma 8.3

$`(X_n)`$ satisfies condition (2) with the coefficients in (3) determined by the equations (9), (8), (10), and $`D=0`$.

Proof. By the Markov property it suffices to show that $`E(X_1^2|X_0,X_2)=Q(X_0,X_2)`$. To this end, as in the proof of Lemma 8.2, we fix a bounded measurable function $`\varphi (X_0,X_2)=_{n,k}\varphi _{n,k}Q_n(X_0)Q_k(X_2)`$. We will verify that

$$E(X_1^2\varphi (X_0,X_2))=E\left(\left(A(X_0^2+X_2^2+BX_0X_2+C)\right)\varphi (X_0,X_2)\right)$$
(33)
by comparing the coefficients in the orthogonal expansions.

From (30) we get $`y^2Q_n^2(y)\mu (dy)=b_{n+1}^2+b_n^2`$, and $`y^2Q_n(y)Q_{n+2}(y)\mu (dy)=b_{n+1}b_{n+2}`$. Moreover, by orthogonality $`y^2Q_n(y)Q_{n+k}(y)\mu (dy)=0`$ except when $`k=0,2,2`$. Using these identities and the expansion $`y^2\varphi (x,z)\nu (dx,dy,dz)=_{n,k}\rho ^{n+k}\varphi _{n,k}y^2Q_n(y)Q_{n+k}(y)\mu (dy)`$ we see that

$`{\displaystyle y^2\varphi (x,z)\nu (dx,dy,dz)}=`$ (34)
$`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\rho ^{2n}\varphi _{n,n}(b_{n+1}^2+b_n^2)+\rho ^2{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\rho ^{2n}b_{n+1}b_{n+2}(\varphi _{n,n+2}+\varphi _{n+2,n}).`$
We now turn to the right hand side of (33). Since

$$x^2\varphi (x,z)\nu (dx,dy,dz)=\varphi (x,z)\underset{n=0}{\overset{\mathrm{}}{}}\rho ^{2n}x^2Q_n(x)Q_n(z)\mu (dx)\mu (dz),$$

a calculation based on (30) yields

$`{\displaystyle x^2\varphi (x,z)\nu (dx,dy,dz)}=`$ (35)
$`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\rho ^{2n}\varphi _{n,n}(b_{n+1}^2+b_n^2)+{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\rho ^{2n}\varphi _{n,n+2}b_{n+1}b_{n+2}+\rho ^4{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\rho ^{2n}\varphi _{n+2,n}b_{n+1}b_{n+2}.`$
By symmetry,

$`{\displaystyle (x^2+z^2)\varphi (x,z)\nu (dx,dy,dz)}=`$ (36)
$`2{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\rho ^{2n}\varphi _{n,n}(b_{n+1}^2+b_n^2)+(1+\rho ^4){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\rho ^{2n}(\varphi _{n,n+2}+\varphi _{n+2,n})b_{n+1}b_{n+2}.`$
Another elementary calculation using (30) gives

$`{\displaystyle xz\varphi (x,z)\nu (dx,dy,dz)}=`$ (37)
$`{\displaystyle \frac{1}{\rho ^2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\rho ^{2n}\varphi _{n,n}b_n^2+\rho ^2{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\rho ^{2n}\varphi _{n,n}b_{n+1}^2+`$
$`\rho ^2{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\rho ^{2n}(\varphi _{n,n+2}+\varphi _{n+2,n})b_{n+1}b_{n+2}.`$
Since $`\varphi (x,z)\nu (dx,dy,dz)=_{n=0}^{\mathrm{}}\rho ^{2n}\varphi _{n,n}`$, we can now verify that (33) holds true by verifying the relation

$$y^2\varphi (x,z)\nu (dx,dy,dz)=\left(A(x^2+z^2)+Bxz+C\right)\varphi (x,z)\nu (dx,dy,dz).$$

Comparing the coefficients in the expansions, the latter reduces to the following.

* Coefficients at $`\varphi _{0,0}`$ match when

$$(2A+B\rho ^2)b_1^2+C=b_1^2$$

(this holds true by (9))
* Coefficients at $`\varphi _{n,n}`$ for $`n1`$ match when

$$2A(b_{n+1}^2+b_n^2)+B\rho ^2b_{n+1}^2+B\frac{b_n^2}{\rho ^2}+C=b_{n+1}^2+b_n^2$$

(this holds true after a longer calculation using $`b_{n+1}^2=qb_n^2+1`$, (8), and (9))
* Coefficients at $`\varphi _{n,n+2}`$ and at $`\varphi _{n+2,n}`$ for $`n0`$ match when

$$A(1+\rho ^4)+B\rho ^2=\rho ^2$$

(this holds true by (8))

This implies that (33) holds true.

## 9 Concluding Remarks

1. To simplify the notation, we restricts ourselves to one-dimensional distributions of the random fields on $``$. In principle our technique is directly applicable to all finite-dimensional distributions, and to multivariate-valued random fields on $``$, see .
2. Symmetric distributions in Theorem 3.2 differ from the distributions corresponding to quadratic conditional variances in .
3. The density in Theorem 3.2 (i) has explicit product representation which can be recovered by a change of variable from the density identified as the “weight function” for the continuous $`q`$-Hermite polynomials, see \[10, Section 3.29\] and \[2, (2.14)\].

An interesting explicit case arises when $`Q(x,y)=\frac{\rho ^2}{1+\rho ^2}\left(x^2+y^2\right)+\frac{\rho ^2(1\rho ^2)}{1+\rho ^2}xy+\frac{(1\rho ^2)(1\rho ^3)}{1+\rho ^2}`$. In this case $`q=0`$ and thus $`Q_n(x/2)/2`$ are the Chebyshev polynomials of the second kind. This implies that $`|X_0|2`$ has the density $`\frac{1}{2\pi }\sqrt{4x^2}`$.
4. It would be interesting to know whether there are stationary processes which satisfy conditions (1), (2) with coefficients in (3) corresponding to $`R>2`$. (Formula \[8, (3.13)\] with $`q^{}=1/q<1`$ indicates that the conclusion of Lemma 8.1 fails.)

### Acknowledgement

I would like to thank M. E. H. Ismail, W. Matysiak, P. Szabłowski, and the referee for helpful comments and information.

#### Note Added Late

Matysiak  pointed out to us that assumption (6) can be weakened to $`1+r_22\rho ^2>0`$; the latter condition is just the linear independence of $`X_0,X_1,X_2`$.