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# Contents

## 1 Introduction

It is, in general, acknowledged that the theoretical treatment of electromagnetic induction phenomena presents some problems when part of the electrical circuit is moving. Some authors speak of exceptions to the flux rule;111 R. Feynman, R. Leighton and M. Sands The Feynman Lectures on Physics, vol. II, (Addison Wesley, Reading, Ma., 1964 ), pp. 17 - 2,3. others save the flux rule by ad hoc choices of the integration line over which the induced $`emf`$ is calculated. Several attempts to overcome these difficulties have been made; a comprehensive one has been performed by Scanlon, Henriksen and Allen.222 P.J. Scanlon, R.N. Henriksen and J.R. Allen, “Approaches to electromagnetic induction”, Am. J. Phys., 37, (1969), 698 - 708. However, their treatment - as others - fails to recognize that one must distinguish between the velocity of the circuit elements and the velocity of the electrical charges (see section 2 below). Therefore, these authors reestablish the flux rule and, consequently, do not solve the problems posed by its application.

Since 1992, I have been teaching electromagnetism in a course for Mathematics students and I had to deal with the problems outlined above. I have found that it is possible to get a general law for electromagnetic induction that contains the standard flux rule as a particular case.

The matter has conceptual relevance; it has also historical and epistemological aspects that deserve to be investigated. Therefore, it is, perhaps, worthwhile to submit the following considerations to the attention of a public wider than that of my students.

## 2  A general law for electromagnetic induction

Textbooks show a great variety of positions about how the flux rule can be applied to the known experimental phenomena of electromagnetic induction. Among the more lucid approaches, let us refer to the treatment given by Feynman, Leighton and Sands in the Feynman Lectures on Physics. They write:

> In general, the force per unit charge is $`\stackrel{}{F}/q=\stackrel{}{E}+\stackrel{}{v}\times \stackrel{}{B}`$. In moving wires there is the force from the second term. Also, there is an $`\stackrel{}{E}`$ field if there is somewhere a changing magnetic field. They are two independent effects, but the emf around the loop of wire is always equal to the rate of change of magnetic flux through it.333 R. Feynman, R. Leighton and M. Sands The Feynman Lectures on Physics, vol. II, (Addison Wesley, Reading, Ma., 1964), p. 17 - 2.

This sentence is followed by a paragraph entitled Exceptions to the “flux rule”, where the authors treat two cases - the Faraday disc and the ‘rocking plates’ - both characterized by the fact that there is a part of the circuit in which the material of the circuit is changing. As the authors put it, at the end of the discussion:444 Ibidem, p. 17 - 3.

> The ‘flux rule does not work in this case. It must be applied to circuits in which the material of the circuit remain the same. When the material of the circuit is changing, we must return to the basic laws. The correct physics is always given by the two basic laws
>
> $`\stackrel{}{F}`$ $`=`$ $`e(\stackrel{}{E}+\stackrel{}{v}\times \stackrel{}{B})`$ (1)
> $`\times \stackrel{}{E}`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{B}}{t}}`$ (2)

### 2.1  A definition of $`emf`$

In order to try shedding some more light on the subject, let us begin with the acknowledgement that the expression of Lorentz force

$$\stackrel{}{F}=q(\stackrel{}{E}+\stackrel{}{v}\times \stackrel{}{B})$$
(3)
not only gives meaning to the fields solutions of Maxwell equations when applied to point charges, but yields new predictions.555 The fact that the expression of Lorentz force can be derived by considering an inertial frame in which the charge is at rest and by assuming that the force acting on it is simply given by $`\stackrel{}{F}^{}=q\stackrel{}{E}^{}`$, does not change the matter.

The velocity appearing in the expression of Lorentz force is the velocity of the charge: from now on, we shall use the symbol $`\stackrel{}{v}_{charge}`$ for distinguishing the charge velocity from the velocity $`\stackrel{}{v}_{line}`$ of the circuit element that contains the charge. This is a basic point of the present treatment.

Let us consider the integral of $`(\stackrel{}{E}+\stackrel{}{v}_{charge}\times \stackrel{}{B})`$ over a closed loop:

$$=_l(\stackrel{}{E}+\stackrel{}{v}_{charge}\times \stackrel{}{B})\stackrel{}{dl}=_l\stackrel{}{E}\stackrel{}{dl}+_l(\stackrel{}{v}_{charge}\times \stackrel{}{B})\stackrel{}{dl}$$
(4)
This integral yields the work done by the electromagnetic field on a unit positive point charge along the closed path considered. It presents itself as the natural definition of the electromotive force, within the Maxwell - Lorentz theory: $`emf=`$.

Let us now calculate the value of $``$ given by equation (4). The calculation of the first integral appearing in the third member of equation (4) yields:

$$_l\stackrel{}{E}\stackrel{}{dl}=_Srot\stackrel{}{E}\widehat{n}𝑑S=_S\frac{\stackrel{}{B}}{t}\widehat{n}𝑑S$$
(5)
where $`S`$ is any surface having the line $`l`$ as contour and where we have made use of Maxwell equation (2). The calculation of the last integral of equation (5) yields:

$$_l\stackrel{}{E}\stackrel{}{d}l=\frac{d}{dt}_S\stackrel{}{B}\widehat{n}𝑑S_l(\stackrel{}{v}_{line}\times \stackrel{}{B})\stackrel{}{d}l$$
(6)
where $`\stackrel{}{v}_{line}`$ is the velocity of the circuit element $`\stackrel{}{dl}`$.666 A. Sommerfeld, Lectures in Theoretical Physics, vol. II, (Academic Press, New York) 1950, pp. 130 - 132; ibidem, vol. III, p. 286. Notice that equation (6) is the result of a theorem of vectorial calculus that is valid for any vector field $`\stackrel{}{G}`$, like the magnetic field, for which $`div\stackrel{}{G}=0`$. Therefore, we get:

$$=\left[\frac{d}{dt}_S\stackrel{}{B}\widehat{n}𝑑S_l(\stackrel{}{v}_{line}\times \stackrel{}{B})\stackrel{}{d}l\right]+_l(\stackrel{}{v}_{charge}\times \stackrel{}{B})\stackrel{}{dl}$$
(7)
This equation says that:

1. The induced emf is, in general, given by three terms.
2. The first two - grouped under square brackets for underlining their common mathematical and physical origin - come from the line integral $`_l\stackrel{}{E}\stackrel{}{dl}`$, whose value is controlled by Maxwell equation (2) through equation (5). Accordingly, their sum must be zero when the magnetic field does not depend on time. In this case, the law assumes the simple form:

$$=_l(\stackrel{}{v}_{charge}\times \stackrel{}{B})\stackrel{}{dl}$$
(8)
3. The third term comes from the magnetic component of Lorentz force; we shall see later how this term may be different from zero.
4. The flux rule is contained in the general law as a particular case.

The general law (7) can be written also in terms of the vector potential $`\stackrel{}{A}`$. If we put

$$\stackrel{}{E}=grad\phi \frac{\stackrel{}{A}}{t}$$
(9)
in the first integral of the third member of equation (4), we get at once

$$=_l\frac{\stackrel{}{A}}{t}\stackrel{}{d}l+_l(\stackrel{}{v}_{charge}\times \stackrel{}{B})\stackrel{}{dl}$$
(10)
since $`_lgra𝑑\phi \stackrel{}{d}l=0`$. When the magnetic field does not depend on time ($`\stackrel{}{B}/t=0`$), we have:

$$_l\stackrel{}{E}\stackrel{}{d}l=_l\frac{\stackrel{}{A}}{t}\stackrel{}{d}l=_Srot\frac{\stackrel{}{A}}{t}\widehat{n}𝑑S=0$$
(11)
since $`rot(\stackrel{}{A}/t)=0`$, because $`\stackrel{}{B}/t=0`$. Therefore, the only surviving term in the expression of the induced $`emf`$ is coming from the magnetic component of Lorentz force and the law assumes the form given by equation (8).

### 2.2 $`\stackrel{}{A}`$ versus $`\stackrel{}{B}`$

Traditionally, textbooks use equations containing the magnetic field in dealing with induction phenomena: the $`emf`$ in a loop is partially (general law) or totally (flux rule) dependent on a surface integral of $`\stackrel{}{B}`$. The laws are predictive, but they are not causal laws. The values of $`\stackrel{}{B}`$ over the surface of integration cannot be causally related to the value of the $`emf`$ around the circuit because $`\stackrel{}{B}`$ acts at a distance (with an infinite propagation velocity): in this case, $`\stackrel{}{B}`$ is not a good field, since a good field, in Feynman’s words, can be defined as

> …a set of numbers we specify in such a way that what happens at a point depends only on the numbers at that point. We do not need to know any more about what’s going on at other places.777 R. Feynman, R. Leighton and M. Sands The Feynman Lectures on Physics, vol. II, (Addison Wesley, Reading, Ma., 1964), p. 15 - 7. Feynman, uses the term ‘real field’ instead of ‘good field’. Feynman deals with this problem in discussing the Bohm - Aharanov effect. It is interesting to see that an identical situation arises in classical electromagnetism.

On the contrary, the same equations written in terms of the vector potential, are causal laws since they relate the value of the $`emf`$ around the circuit to the values of $`\stackrel{}{A}/t`$ at the points of the loop: the vector potential is, in this case, a good field.

## 3 How the law works

### 3.1 General features

Let us come back to the general law (7 or 10). The charge velocity appearing in these equations contributes to build up, through the factor $`\stackrel{}{v}_{charge}\times \stackrel{}{B}`$ the induced electromotive field. The charge velocity is given by:

$$\stackrel{}{v}_{charge}=\stackrel{}{v}_{line}+\stackrel{}{v}_{drift}$$
(12)
where $`\stackrel{}{v}_{drift}`$ is the drift velocity.

Therefore, the general equation for electromagnetic induction assumes the form:

$``$ $`=`$ $`\left[{\displaystyle \frac{d}{dt}}{\displaystyle _S}\stackrel{}{B}\widehat{n}𝑑S{\displaystyle _l}(\stackrel{}{v}_{line}\times \stackrel{}{B})\stackrel{}{d}l\right]+`$ (13)
$`+`$ $`{\displaystyle _l}(\stackrel{}{v}_{line}\times \stackrel{}{B})\stackrel{}{dl}+{\displaystyle _l}(\stackrel{}{v}_{drift}\times \stackrel{}{B})\stackrel{}{dl}`$
or, in terms of the vector potential:

$$=_l\frac{\stackrel{}{A}}{t}\stackrel{}{d}l+_l(\stackrel{}{v}_{line}\times \stackrel{}{B})\stackrel{}{dl}+_l(\stackrel{}{v}_{drift}\times \stackrel{}{B})\stackrel{}{dl}$$
(14)
When the circuit is made by a loop of wire, equation (13) reduces to equation

$$=\left[\frac{d}{dt}_S\stackrel{}{B}\widehat{n}𝑑S_l(\stackrel{}{v}_{line}\times \stackrel{}{B})\stackrel{}{d}l\right]+_l(\stackrel{}{v}_{charge/line}\times \stackrel{}{B})\stackrel{}{dl}$$
(15)
and equation (14) to equation

$$=_l\frac{\stackrel{}{A}}{t}\stackrel{}{d}l+_l(\stackrel{}{v}_{charge/line}\times \stackrel{}{B})\stackrel{}{dl}$$
(16)
because the drift velocity is always parallel to the line element $`\stackrel{}{dl}`$ and, consequently, the integral containing it is zero. We have introduced the new notation $`v_{charge/line}`$ for remembering that we are dealing with the velocity of the charge that, in this case, can be replaced by the velocity of the circuit element that contains it. It is worth stressing again that, when the magnetic field does not depend on time, equation (15) cannot be read in terms of flux variation, since the sum of the first two terms under square brackets is zero as it is the equivalent term of equation (16) containing the time derivative of the vector potential.

When part of the circuit is made by extended material, the calculation of the drift velocity contribution to the induced $`emf`$ is not an easy task, since the distribution of current lines in the material may be very complicated.888 The distribution of currents in extended materials have been studied since the mid of eighteenth century. The following references are those that I know: E. Jochmann, “On the electric currents induced by a magnet in a rotating conductor”, Phil. Mag., 27, (1864), 506 - 528; Phil. Mag., 28, (1864), 347 - 349; H. Hertz, ‘On induction in rotating spheres’, in Miscellaneous papers, (J.A. Barth, Leipzig) 1895, pp. 35 - 126. A discussion of some aspects of Hertz’s work can be found in: Jed Z. Buchwald, The creation of scientific effects - Heinrich Hertz and electric waves, (The University of Chicago Press, Chicago and London) 1994, pp. 95 - 103; the papers by Boltzmann and Corbino quoted in footnotes 12; 13. V. Volterra, “Sulle correnti elettriche in una lamina metallica sotto l’azione di un campo magnetico”, Il Nuovo Cimento, 9, 23 - 79 (1915). One should also see the literature on eddy currents.  In section 3.2.6 and 3.2.7 we shall treat particularly simple cases with a circular symmetry. In all other cases, we shall neglect the drift velocity in the calculation of the induced $`emf`$ (look at the end of section 3.2.6 for a further discussion of this point).

### 3.2 How it works in specific cases

We shall now discuss some cases widely treated in literature, in order to see how the general law (and the flux rule) can be applied to specific problems.

#### 3.2.1 When a bar is moving

Let us consider the circuit of fig. 1. The conducting bar $`A`$, of length $`a`$, slides with constant velocity $`\stackrel{}{v}`$ over the conducting frame $`T`$ in a uniform and constant magnetic field $`\stackrel{}{B}`$ perpendicular to the plane of the frame and entering the page: the frame $`T`$ is at rest with respect to the source of the magnetic field. A steady current flows along the direction indicated in the figure. Notice that in this case we can neglect without approximation the drift velocity of the charge because it is always directed along the integration line (owing to the Hall effect). According to the general law (16), the first integral is zero as zero is sum of the first two terms of equation (15): as a matter of fact, if we choose the counterclockwise direction for the line integral, we obtain for the sum of these two terms:

$$Bav_{line}Bav_{line}=0$$
(17)
The induced electromotive force is then given by:

$$=Bav_{charge/line}$$
(18)
The surviving term is the one coming from the magnetic component of Lorentz force. Moreover, the theory predicts that the emf is localized into the bar: the bar acts as a battery and, as a consequence, between the two points $`1`$ and $`2`$ of the frame that are in contact with the bar, one should measure a potential difference given by $`V_1V_2=ir`$, where $`i`$ is the circulating current and $`r`$ is the resistance of the bar.

Finally, it is worth stressing that the energy balance shows that the magnetic component of Lorentz force plays the role of an intermediary: the electrical energy dissipated in the circuit comes from the work done by the force that must be applied to keep the bar moving with constant velocity.

Let us now recall how the flux rule deals with this case. It predicts an emf given by $`Bav_{line}`$. In the light of the general law (15) and of its discussion, we understand why the flux rule predicts correctly the value of the emf: the reason lies in the fact that the two line integrals under square brackets cancel each other. However, we have shown above that the physics embedded in the general equation (15) forbids to read equation (18) as the result of:

$$=Bav_{line}+[Bav_{line}+Bav_{line}]=Bav_{line}$$
(19)
that leaves operative the first term coming from the flux variation.

Before leaving this subject, it is worth saying something about the long debate about the meaning of ‘localized $`emf`$’. With reference to the moving bar, we can proceed as follows. The points $`1`$ and $`2`$ divide the circuit into two parts: the one on the left obeys Ohm law $`\mathrm{\Delta }V=iR`$ and the current enters it at point $`1`$ at higher potential and leaves it at point $`2`$ at lower potential; the part on the right does not obeys Ohm law and the current enters it at point $`2`$ at lower potential and leaves it at point $`1`$ at higher potential. We say that the $`emf`$ is localized in the right part of the circuit: it can be experimentally distinguished from the other.

#### 3.2.2 The Faraday disc: 1

A conducting disc is rotating with constant angular velocity $`\omega `$ in a uniform and constant magnetic field $`\stackrel{}{B}`$ perpendicular to the disc and entering the page (fig. 2). A conducting frame $`DABC`$ makes conducting contacts with the center and a point on the periphery of the disc. A steady current $`i`$ circulates as indicated by the arrow.

By applying the general laws (16) or (15) to the integration line at rest $`ABCDA`$, we see that, if we neglect the contribution of the drift velocity:

* since the magnetic field does not depend on time, the first integral in (16) is zero
* equivalently, the sum of the first two integrals in (15) is zero. In the present case, each of the two terms is zero: the flux associated with the circuit does not change and $`\stackrel{}{v}_{line}`$ is zero everywhere, since the integration line has been chosen at rest
* the only surviving term is the one due to the magnetic component of Lorentz force and its value is given by (we are taking the counterclockwise direction for the line integration):

$$=\frac{1}{2}\omega BR^2$$
(20)
where $`R`$ is the disc radius. In doing this calculation, we have neglected - as explained at the end of section 3.1, the charge drift velocity (we shall take it into account in section 3.2.7).

Of course, the emf given by equation (20) is induced along any radius, as it can be seen by considering the circuit $`ABCC^{}DA`$. If the radius $`C^{}D`$ is considered at rest, the case is the same as that of the circuit $`ABCDA`$ discussed just before. If the radius $`C^{}D`$ is considered in motion with angular velocity $`\omega `$, we see again, both from equation (16) and equation (15), that the only surviving term is the one due to the magnetic component of Lorentz force with the only difference that, now, the first term of equation (15) $`(1/2)\omega BR^2`$ is cancelled out by the second $`(1/2)\omega BR^2`$, due to the movement of the line elements of the radius $`DC^{}`$; the line elements lying on the $`CC^{}`$ arc give a null contribution.

Let us now see how the flux rule deals with the Faraday disc. If we choose the lines at rest $`ABCDA`$ or $`ABCC^{}DA`$, we find that the induced emf is zero; if we choose the integration line $`ABCC^{}DA`$, with the radius $`C^{}D`$ considered in motion, we obtain the correct result $`=(1/2)\omega BR^2`$.

As in the case of the moving bar, the flux rule gets the correct result only because the two line integrals under square brackets in equation (15) cancel each other. As in the case of the moving bar, the physics embedded in the general equation forbids an interpretation of the mathematical result in terms of flux variation: again the physical origin of the induced emf is due to the intermediacy of the magnetic component of Lorentz force.

#### 3.2.3 The unipolar induction

The so called unipolar induction is illustrated in fig. 3. When the cylindrical and conducting magnet $`M`$ rotates about its axis with angular velocity $`\omega `$ in the counterclockwise direction, a current flows in the circuit as indicated by the arrow. It is easy to see that the discussion of the unipolar induction can be reduced to that of the Faraday disc for both general law and flux rule: the general law applied to the circuit $`ABCDEA`$ at rest yields an $`emf`$ induced in the radius $`DE`$ given by $`(1/2)\omega B\overline{DE}^2`$, while the flux rule applied to the same circuit predicts a zero emf. If one consider instead the integration line $`ABCDD^{}EA`$ with the radius $`D^{}E`$ in motion, analogous to the path $`ABCC^{}DA`$ of fig. 2 concerning the Faraday disc, one can follow the same arguments developed in that case.

#### 3.2.4 The flux varies, but…

In fig. 4 a case discussed by Scanlon, Henriksen and Allen999 P.J. Scanlon, R.N. Henriksen and J.R. Allen, “Approaches to electromagnetic induction”, Am. Jour. Phys., 37, pp. 705 - 706 (1969). and originally due to Kaempffer101010 F. A. Kaempffer, Elements of Physics (Blaisdell Publ. Co., Waltham, Mass., 1967), p. 164; quoted by Scanlon, Henriksen and Allen. is presented. When the conducting magnet $`M`$ moves, as indicated, with constant velocity $`\stackrel{}{v}`$, there is no induced $`emf`$ \- if the magnet is sufficiently long in the direction perpendicular to the page so that there is a magnetic field only within the magnet. If we consider the integration line $`ACDEFGA`$ at rest, there is a flux variation without induced $`emf`$. Also in this case, we must choose ad hoc the integration line - for instance the line $`ACC^{}D^{}DEFGA`$, where the segment $`C^{}D^{}`$ moves with the magnet - if we want to save the flux rule. On the other hand, the general law works well, whatever integration line is chosen.

#### 3.2.5 The ‘rocking plates’

This is one of the two ‘exceptions to the flux rule’ discussed by Feynman, Leighton and Sands (fig. 5).111111 R. Feynman, R. Leighton and M. Sands The Feynman Lectures on Physics, (Addison Wesley, Reading, Ma., 1964 ), pp. 17 - 3. The two plates oscillate slowly back and forth so that their point of contact moves from $`P`$ to $`P^{}`$ and viceversa. The circuit is closed by a wire that connects point $`A`$ and $`C`$. The magnetic field $`\stackrel{}{B}`$ is perpendicular to the plates and enters the page. The authors write:

> If we imagine the “circuit” to be completed through the plates on the dotted line shown in the figure, the magnetic flux through this circuit changes by a large amount as the plates are rocked back and forth. Yet the rocking can be done with small motions, so that $`\stackrel{}{v}\times \stackrel{}{B}`$ is very small and there is practically no emf.

According to the general law (15) the sum of the first two terms of equation (15) must be zero. Hence:

$$_l(\stackrel{}{v}_{line}\times \stackrel{}{B})\stackrel{}{d}l=\frac{d}{dt}_S\stackrel{}{B}\widehat{n}𝑑S$$
(21)
Therefore, the third term of equation (15) coming from the magnetic component of Lorentz force equals the magnetic flux variation changed in sign, if, of course, we assume that the charge velocity can be taken equal to the velocity of the element of conductor that contains it (the drift velocity is neglected). The conclusion is that there is an induced $`emf`$: its average value is given by $`\mathrm{\Delta }\varphi /\mathrm{\Delta }t`$ where $`\mathrm{\Delta }\varphi `$ is the flux variation between the two extreme positions of the plates and $`\mathrm{\Delta }t`$ the interval of time taken in going from one position to the other: the induced $`emf`$ gets smaller and smaller as $`\mathrm{\Delta }t`$ gets larger and larger; when the motion of the plates is very slow, we can conclude with Feynman that ‘there is practically no $`emf`$’. Notice that the induced $`emf`$ changes in sign (and the current its direction) when the the rocking motion is reversed.

#### 3.2.6 The Corbino disc

The discussion of this case will show how the charge drift velocity, always neglected before, plays its role in the building up of the induced electromotive field and, therefore, of the induced $`emf`$. In 1911, Corbino studied theoretically and experimentally the case of a conducting disc with a hole at its center. If a voltage is applied between the inner and the outer periphery of the disc, a radial current will flow, provided that the experimental setup is realized in a way suitable for maintaining the circular symmetry: the inner and outer periphery are covered by highly conducting electrodes; therefore, the inner and outer periphery are two equipotential lines. If a uniform and constant magnetic field is applied perpendicularly to the disc, a circular current will flow in the disc.121212 O.M. Corbino, “Azioni elettromagnetiche dovute agli ioni dei metalli deviati dalla traiettoria normale per effetto di un campo”, Il Nuovo Cimento 1, 397 - 419 (1911). A german translation of this paper appeared in Phys. Zeits., 12, 561 - 568 (1911). For a historical reconstruction see: S. Galdabini and G. Giuliani, “Magnetic field effects and dualistic theory of metallic conduction in Italy (1911 - 1926): cultural heritage, epistemological beliefs, and national scientific community”, Ann. Science 48, 21 - 37 (1991). As pointed out by von Klitzing, the quantum Hall effect may be considered as an ideal (and quantized) version of the Corbino effect corresponding to the case in which the current in the disc, with an applied radial voltage, is only circular: K. von Klitzing, “The ideal Corbino effect”, in: P.E. Giua ed., Commemorazione di Orso Mario Corbino, (Centro Stampa De Vittoria, Roma, 1987), pp. 43 - 58.

The first theoretical treatment of this case is due, as far as I know, to Boltzmann who wrote down the equations of motion of charges in combined electric and magnetic fields.131313 L. Boltzmann, Anzeiger der Kaiserlichen Akademie der Wissenschaften in Wien,23, (1886), 77 - 80; Phil. Mag., 22, (1886), 226 - 228. Corbino, apparently not aware of this fact, obtained the same equations already developed by Boltzmann. However, while Boltzmann focused on magnetoresistance effects, Corbino interpreted the theoretical results in terms of radial and circular currents and studied experimentally the magnetic effects due to the latter ones.

The application of the general law of electromagnetic induction to this case leads to the same results usually obtained (as Boltzmann and Corbino did) by writing down and solving the equations of motion of the charges in an electromagnetic field (by taking into account, explicitly or implicitly, the scattering processes).

If $`I_{radial}`$ is the radial current, the radial current density $`J(r)`$ will be:

$$J(r)=\frac{I_{radial}}{2\pi rs}$$
(22)
and the radial drift velocity:

$$v(r)_{drift}=\frac{I_{radial}}{2\pi rsne}$$
(23)
where $`s`$ is the thickness of the disc, $`n`$ the electron concentration and $`e`$ the electron charge. In the present case the general law of electromagnetic induction assumes the simple form of equation (8) with $`v_{charge}=v_{drift}`$; therefore, the induced $`emf`$ around a circle of radius $`r`$ is given by:

$$_{circular}=_0^{2\pi r}(\stackrel{}{v}(r)_{drift}\times \stackrel{}{B})\stackrel{}{dr}=\frac{I_{radial}B}{sne}$$
(24)
The circular current $`dI(r)_{circular}`$ flowing in a circular strip of radius $`r`$ and section $`sdr`$ will be, if $`\rho `$ is the resistivity:

$$dI_{circular}=\frac{_{circular}sdr}{\rho \mathrm{\hspace{0.17em}2}\pi r}=\frac{\mu B}{2\pi }I_{radial}\frac{dr}{r}$$
(25)
and the total circular current:

$$I_{circular}=\frac{\mu B}{2\pi }I_{radial}\mathrm{ln}\frac{r_2}{r_1}$$
(26)
where $`\mu `$ is the electron mobility, $`r_1`$ and $`r_2`$ the inner and outer radius of the disc (we have used the relation $`\mu =1/\rho ne`$). Equation (26) is the same as that derived and experimentally tested by Corbino.

The power dissipated in the disc is:

$$W=(I^2R)_{radial}+(I^2R)_{circular}=I_{radial}^2R_{radial}(1+\mu ^2B^2)$$
(27)
where we have used equation (26) and the two relations:

$`R_{radial}`$ $`=`$ $`{\displaystyle \frac{\rho }{2\pi s}}\mathrm{ln}{\displaystyle \frac{r_2}{r_1}}`$ (28)
$`R_{circular}`$ $`=`$ $`{\displaystyle \frac{\rho ^2}{s^2}}{\displaystyle \frac{1}{R_{radial}}}`$ (29)
Equation (27) shows that the phenomenon may be described as due to an increased resistance $`R_{radial}(1+\mu ^2B^2)`$: this is the magnetoresistance effect. The circular induced $`emf`$ is ‘distributed’ homogeneously along each circle. Each circular strip of section $`sdr`$ acts as a battery that produces current in its own resistance: therefore, the potential difference between two points arbitrarily chosen on a circle is zero. Hence, as it must be, each circle is an equipotential line.

The above application of the general law yields a description of Corbino disc that combine Boltzmann (magnetoresistance effects) and Corbino (circular currents) point of view and shows how the general law can be applied to phenomena traditionally considered outside the phenomenological domain for which it has been derived.

#### 3.2.7 The Faraday disc: 2

The discussion of Corbino disc helps us in better understanding the physics of the Faraday disc. Let us consider a Faraday disc in which the circular symmetry is conserved. This may be difficult to realize; anyway, it is interesting to discuss it. As shown above, the steady condition will be characterized by the flow of a radial and of a circular current. The mechanical power needed to keep the disc rotating with constant angular velocity $`\omega `$ is equal to the work per unit time done by the magnetic field on the rotating radial currents. Then, it will be given by:

$$W=_0^{2\pi }_{r_1}^{r_2}(J_{radial}rd\alpha s)(Bdr)(\omega r)=I_{radial}\frac{1}{2}\omega B(r_2^2r_1^2)$$
(30)
where the symbols are the same as those used in the previous section. The point is that the term

$$=\frac{1}{2}\omega B(r_2^2r_1^2)$$
(31)
is the induced $`emf`$ due only to the motion of the disc. This $`emf`$ is the primary source of the induced currents, radial and circular. Therefore, the physics of the Faraday disc with circular symmetry, may be summarized as follows:

1. the source of the induced currents is the induced $`emf`$ due to the rotation of the disc
2. the primary product of the induced $`emf`$ is a radial current
3. the drift velocity of the radial current produces in turn a circular induced $`emf`$ that give rise to the circular current

## 4 General law and flux rule

The cases discussed above must be considered as illustrative examples. Obviously, the conceptual foundations of the general law and of the flux rule can be discussed without any reference to particular cases. The general law has been already dealt with in great detail. Let us now focus our attention on the flux rule.

Let us define - this, of course, is not our choice - the induced emf as

$$=_l\stackrel{}{E}\stackrel{}{d}l$$
(32)
where $`\stackrel{}{E}`$ obeys Maxwell equations. As shown above (equations 5 \- 6), this leads immediately to the conclusion that the induced emf is given by:

$$=_l\stackrel{}{E}\stackrel{}{d}l=\frac{d}{dt}_S\stackrel{}{B}\widehat{n}𝑑S_l(\stackrel{}{v}_{line}\times \stackrel{}{B})\stackrel{}{d}l$$
(33)
or, in terms of the vector potential, by (equations 10 \- 11):

$$=_l\stackrel{}{E}\stackrel{}{d}l=_l\frac{\stackrel{}{A}}{t}\stackrel{}{d}l$$
(34)
and is always zero when $`\stackrel{}{B}/t=0`$.

Of course, this result is not satisfactory, if we want to save the flux rule. Therefore, one may suggest to define the induced emf as:

$$=_l\stackrel{}{E}\stackrel{}{d}l+_l(\stackrel{}{v}_{line}\times \stackrel{}{B})\stackrel{}{d}l=\frac{d}{dt}_S\stackrel{}{B}\widehat{n}𝑑S$$
(35)
thus reestablishing the flux rule. This is, for instance, the choice made by Scanlon, Henriksen and Allen.141414 P.J. Scanlon, R.N. Henriksen and J.R. Allen, “Approaches to electromagnetic induction”, Am. Jour. Phys., 37, (1969), p. 701.  However:

* since the velocity appearing in equation (35) is clearly the velocity of the line element $`\stackrel{}{d}l`$, this definition has no physical meaning.
* the definition is ad hoc. As a consequence, as in the paper by Scanlon, Henriksen and Allen, the line of integration must be chosen ad hoc for obtaining predictions in accordance with the experimental findings.

## 5 Conclusions

The straightforward application of two fundamental laws of electromagnetism – Maxwell equation $`rot\stackrel{}{E}=\stackrel{}{B}/t`$ and Lorentz law $`\stackrel{}{F}=q(\stackrel{}{E}+\stackrel{}{v}\times \stackrel{}{B})`$ – leads to a general law for electromagnetic induction phenomena: it includes, as a particular case, the standard flux rule. The treatment given in this paper shows that standard derivations:

1. Fail to recognize that the velocity involved in the integration of Maxwell equation (2) is the velocity of the line elements composing the closed path of integration and not the velocity of the charges.
2. Overlook the fact that the expression of Lorentz force constitutes a fundamental postulate not included in Maxwell equations; this postulate not only gives physical meaning to the fields solutions of Maxwell equations – when applied to charges – but allows also the prediction of new phenomena.
3. In relevant cases, when a part of the circuit is moving, they are misleading as far as the physical origin of the phenomenon is concerned: they attribute the effect to a flux variation when, instead, the origin lies in the intermediacy of the magnetic component of Lorentz force.

Finally, the present derivation:

* provides with a rigorous mathematical treatment the statement according to which electromagnetic induction phenomena are the product of two independent processes: time variation of magnetic field and effects of magnetic field on moving charges.
* yields a law that can be successfully applied to phenomena - those presented by Corbino disc - till now considered outside the phenomenological domain of electromagnetic induction.

Acknowledgements. The author would like to thank Ilaria Bonizzoni for a critical reading of the manuscript and Giuliano Bellodi and Gianluca Introzzi for valuable discussions.


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# Making classical and quantum canonical general relativity computable through a power series expansion in the inverse cosmological constant

## Abstract

We consider general relativity with a cosmological constant as a perturbative expansion around a completely solvable diffeomorphism invariant field theory. This theory is the $`\mathrm{\Lambda }\mathrm{}`$ limit of general relativity. This allows an explicit perturbative computational setup in which the quantum states of the theory and the classical observables can be explicitly computed. The zeroth order corresponds to highly degenerate space-times with vanishing volume. Perturbations give rise to space-times with non-vanishing volumes in a natural way. The spectrum of area- and volume-related observables constructed by coupling the theory to matter can be directly assessed. An unexpected relationship arises at a quantum level between the discrete spectrum of the volume operator and the allowed values of the cosmological constant.

The recent years have been very fertile for the development of canonical quantizations of general relativity. There are now available more than one  highly non-trivial mathematically consistent, finite, anomaly-free theories of canonical quantum gravity. Instrumental in these developments have been the underlying advances in mathematical techniques for dealing with infinite dimensional nonlinear spaces, like the theory of cylindrical functions and associated measures , and the introduction of spin networks to eliminate the over-completeness of the Wilson loop basis  (for a recent summary, see the review article by Rovelli ). In spite of the advances, most of the results obtained from these theories up to now concern statements made at a kinematical level (without imposing the Hamiltonian constraint). One of the remaining main challenges is to obtain results that hold at a dynamical level, imposing the full set of constraints of the theory. This requires making the theories computable, in the sense of obtaining classical observables to promote to quantum observables and solving for the physical set of states that are annihilated by all the constraints. Only in this setting will one be able to make contact with trustworthy physical predictions.

The fact that implementing the full dynamics of canonical quantum gravity is hard to do has plagued the subject since its inception in the 60’s by DeWitt. For instance, one does not know even at a classical level how to compute observable quantities for canonical general relativity (quantities that are invariants under the symmetries of the theory or —in the canonical language— that commute with the constraints), and it is known that such quantities are potentially quite involved . As a consequence, it is very hard to make contact with the semiclassical limit, since one does not have at hand observable quantities to compute. Moreover, the space of quantum states that are annihilated by the quantum constraints cannot be straightforwardly computed . None of these issues was directly tackled by the significant amount of progress that was made in terms of the introduction of connection type variables and loop representations for quantum gravity that we mentioned at the beginning. There is general consensus that an approximation scheme is desirable in such a way as to make the above problems solvable. The challenge consists in finding an approximation scheme that will not conflict with non-perturbative nature of the canonical quantization, in particular the diffeomorphism invariance of the theory.

In this paper we would like to propose a computational scheme that allows to explicitly address the issue of finding observables for the theory (at a classical level) and to find quantum states annihilated by all the constraints. The central idea will be to consider general relativity as a power series expansion in terms of the inverse cosmological constant, and to view the full theory as a “perturbation” of a zeroth order solvable theory. The zeroth order theory corresponds to general relativity in the limit in which the cosmological constant $`\mathrm{\Lambda }\mathrm{}`$. The use of an approximation based on a large cosmological constant has been pioneered in a different context (that of studying the spectrum of the physical Hamiltonian obtained by “de-parameterizing” the theory) by Smolin  and also by Rovelli , with some level of success. We will however be able to make further progress, in significant part due to the new maturity of the field we mentioned above but also because we push the perturbative approach at all levels: constraints, observables and equations of motion, both classically and quantum mechanically. It might be argued that a scheme where the cosmological constant is large (in Planck units) lacks physical relevance, given is small present value. However, having a regime of the theory in which it is completely solvable can be a valuable tool, even if such regime is unphysical. Moreover, as we will discuss below, there is some evidence that our approach might be useful beyond its apparent domain of validity. It might be the case that the situation is analogous to the large $`N`$ approximation in QCD or the strong coupling expansions of lattice field theory, that a priori look less promising that what they end up being.

The starting point is to consider canonical quantum gravity with a cosmological constant. The Hamiltonian constraint of the theory is,

$$H(N)=H_0(N)+\mathrm{\Lambda }d^3xN(x)\sqrt{\mathrm{det}q(x)}$$
(1)
where $`N(x)`$ is a smearing function, $`H_0(N)`$ is the Hamiltonian constraint without a cosmological constant and $`\mathrm{det}q`$ is the determinant of the spatial metric. If one now considers the limit $`\mathrm{\Lambda }\mathrm{}`$, and re-scales the constraint by $`1/\mathrm{\Lambda }`$  one is left with a theory, which we will call “zeroth order” theory for which the Hamiltonian constraint is, $`H^{(0)}(N)=d^3xN(x)\sqrt{\mathrm{det}q(x)},`$ that is, the Hamiltonian constraint is just the square root of the determinant of the spatial metric. In addition to this the theory has the ordinary diffeomorphism constraint and if one uses Ashtekar variables there will also be a Gauss law constraint (both are independent of $`\mathrm{\Lambda }`$). Imposing classically the Hamiltonian constraint of the zeroth order theory, one immediately sees that it corresponds to metrics of identically vanishing determinant, that is, degenerate metrics.

The dynamics of the zeroth order theory is easily solved. One can readily construct quantities that have vanishing Poisson bracket with the Hamiltonian constraint of the zeroth order theory. For instance, one can consider any function depending only on the three metric (or if one is using Ashtekar variables $`(\stackrel{~}{E}_i^a,A_b^j)`$, the densitized triads $`\stackrel{~}{E}_i^a`$) and not on its canonically conjugate momenta. To be a genuine observable the quantity should also have vanishing Poisson bracket with the diffeomorphism constraint. A way of constructing quantities with this property is to couple the theory with matter fields and physically define geometric objects like surfaces or volumetric regions through properties of the matter fields  and then compute the areas or volumes of these regions. The usual problem with this construction is that the addition of matter adds terms to the Hamiltonian constraint as well, and the areas and volumes computed as we just described do not commute with the Hamiltonian constraint. However, in the $`\mathrm{\Lambda }\mathrm{}`$ limit the matter terms drop out of the Hamiltonian and one ends up (in the zeroth order theory) with genuine observables (this is similar to what happens in models without a Hamiltonian constraint as was discussed by Rovelli in the Husain–Kuchař  model ). A concrete example of an observable for the zeroth order theory is obtained by coupling the theory to an antisymmetric tensor field with the standard Kalb–Ramond free Hamiltonian as proposed by Smolin . The antisymmetric tensor field is naturally associated with surfaces of which one can compute the area. The resulting quantity, which essentially is a contraction of the triads with the canonically conjugate momenta of of the Kalb–Ramond field, is naturally diffeomorphism invariant. One can build a “loop-surface” representation (see  for details) for the coupled theory by expanding the gauge invariant wavefunctions in a basis of states $`|\mathrm{\Sigma },s>`$ labeled by equivalence classes under diffeomorphisms of spin networks $`s`$ and three-surfaces $`\mathrm{\Sigma }`$. A natural generalization of the Ashtekar-Lewandowski  measure allows to introduce an inner product. On this basis the area defined by the fields has a discrete spectrum.

The Hamiltonian constraint of the zeroth order theory can be readily realized in the kinematical space  of non-diffeomorphism invariant wavefunctions as,

$$H^{(0)}(N)|\mathrm{\Sigma },s>=\underset{ϵ0}{lim}\underset{vs}{}N(v)V(v,ϵ)|\mathrm{\Sigma },s>$$
(2)
where we have represented, following Thiemann , the determinant of the metric as the volume operator of an infinitesimal region of size $`ϵ`$ centered on the vertices $`v`$ of the spin network. The operator is independent on $`ϵ`$ on this space of states and therefore the above limit is well defined. To consider the action on diffeomorphism invariant states one evaluates the inner product $`<\{\mathrm{\Sigma }^{}\},\{s\}|H^{(0)}(N)|\mathrm{\Sigma },s>`$ which is well defined. We can readily find the space of states that are annihilated by the Hamiltonian constraint of the zeroth order theory: they correspond to states based on spin networks such that the localized volume operator vanishes on each of their vertices. For trivalent vertices this is the case for all states. For four and higher valent vertices the operator vanishes automatically if the vertices are planar (this is not true for the volume operator of , but it holds for the definition of ). If the vertices are not planar it is also possible to find values of the valences of the incoming edges such that the operator vanishes. We collectively denote all these states as $`|\mathrm{\Sigma },s^0>`$. Since the space of solutions of the Hamiltonian constraint is a subspace of the states we were considering, it can be readily endowed with an inner product which is just the restriction to this space of the inner product derived from the Ashtekar–Lewandowski measure. Notice that the area observables we introduced are well defined on this space.

We have therefore seen that the zeroth order theory is well defined. It is non-perturbatively quantized in a basis of spin networks and surfaces, one can solve for the physical states, construct observables, and introduce a physical inner product. This is therefore a sound starting point for a perturbative expansion. The resulting perturbative theory is remarkably similar to bound state perturbation theory in quantum mechanics in the degenerate case. The operator $`H^{(0)}`$ (which can be readily diagonalized in the basis of states considered) has a degenerate spectrum.

To consider the higher order corrections to the zeroth order theory, one is interested in solving the eigenvalue problem

$$<\varphi (\mathrm{\Lambda }^1)|H^{(0)}+\mathrm{\Lambda }^1H^{(1)}=ϵ(\mathrm{\Lambda })<\varphi (\mathrm{\Lambda }^1|$$
(3)
with eigenvalue $`ϵ(\mathrm{\Lambda })`$ equal to zero up to the desired order in perturbation theory. We will assume that the states and the eigenvalue can be expanded in power series in the inverse cosmological constant, and address the resulting equations order by order. Since the spectrum of the zeroth order Hamiltonian is degenerate, we apply degenerate perturbation theory to solve this problem. One starts by considering a state with eigenvalue $`ϵ^{(0)}`$ for the zeroth order Hamiltonian constraint, that is, a spin network state with a definite volume. Notice that this is not a state of the zeroth order theory. Considering first order corrections and solving the equations resulting from (3) one determines the first order corrections and a polynomial $`ϵ(\mathrm{\Lambda }^1)`$ which, set to zero, determines the value of $`\mathrm{\Lambda }`$ for which the perturbative treatment is valid. That is, given an eigenvalue of the zeroth order Hamiltonian, one constructs a solution of the first order theory that has vanishing Hamiltonian for a particular value of lambda. That is, as a consequence of the discrete spectrum of the zeroth order Hamiltonian inherited from the discreteness of the volume operator, the cosmological constant is, in this approach, quantized.

This might sound unusual, but imposing the constraint perturbatively is analogous to considering an atom in a magnetic field and trying to find the states with a given energy level. Such energy level can be reached by choosing a nearby energy state of an unperturbed atom and then tuning the magnetic field to a given value. Various states could correspond to the same energy, depending on the value of the magnetic field. The approximate value of the magnetic field required will also depend on how high one goes in perturbation theory. A detailed discussion of how this is implemented in several systems can be seen in .

A remarkable aspect is that all the calculations involved in the perturbative approach are completely explicit when one works in the spin network basis. One starts by considering eigenvalues and eigenvectors of the zeroth order Hamiltonian. Since this Hamiltonian is the volume operator, which is well understood, the space of states is under control (even when it is non-vanishing). The equation that determines the first order correction involves evaluating the Hamiltonian constraint on the states with a given volume. Such calculations are completely explicit given, for instance, the Hamiltonian introduced by Thiemann , or the Hamiltonian in the Vassiliev arena . The only subtle element in the calculation is that given the infinite degeneracy of the volume operator, the first order correction might involve infinite superpositions and therefore may lead us to corrections to the kinematical inner product. It is worthwhile noticing that in general the first order corrections have non-zero volume, even if one started from the zeroth order state with vanishing volume, as long as one did not restrict to either planar or trivalent spin networks. So the resulting theory is a theory of non-vanishing volume space-times.

The above construction has a classical counterpart when one tries to find observables for the theory. An observable for the first order theory can be written as a perturbation of the observables of the zeroth order level theory we discussed above,

$$O_\mathrm{\Lambda }(\pi ,\stackrel{~}{E})=O^{(0)}(\pi ,\stackrel{~}{E})+\mathrm{\Lambda }^1O^{(1)}(\pi ,\stackrel{~}{E}).$$
(4)
One would like to request that these observables have vanishing Poisson brackets with the Hamiltonian of the theory (more precisely, weakly vanishing, but the extension to such case is immediate). Expanding such requirement in powers of $`\mathrm{\Lambda }^1`$ one gets to zeroth and first order in $`\mathrm{\Lambda }^1`$,

$`\{O^{(0)},H^{(0)}\}`$ $`=`$ $`0`$ (5)
$`\{O^{(1)},H^{(0)}\}+\{O^{(0)},H^{(1)}\}`$ $`=`$ $`0`$ (6)
The first equation determines $`O^{(0)}`$ and the second one leads to a (functional) linear partial differential equation for $`O^{(1)}`$. The construction can be readily continued to higher orders. In all cases one obtains a linear partial differential equation, albeit with a more and more complex inhomogeneous term. It should be noted that one can obtain many observables starting with different solutions to the first equation. The linear partial differential equations are not hard to solve, given that the coefficients of the derivatives are functions of the triads, whereas the derivatives are with respect to the connections. Given its simplicity, the system is always integrable and therefore yields all the observables of the theory. The solutions to these equations will generically be ill defined or non-convergent in certain regions of phase space. This is expected in theories with complicated dynamics and is not a pathology of the method per se but rather a feature of the observable quantities. In particular, one might be concerned that the current approach produces too many observables. For instance, there exist pathological systems that have less observables than degrees of freedom (e.g. the heavy asymmetric top). A priori our method seems to yield in these cases extra observables, but on closer examination the quantities constructed are not acceptable (in the case of the top they take values on a restricted portion of phase space not preserved by evolution). We discuss several pathological examples and how the method correctly handles them in .

Let us consider, as an simple example of interest in gravity, the computation of an observable for Bianchi cosmologies. Observables for these models (with cosmological constant) had never been studied before. We will write the Bianchi models in terms of Ashtekar’s variables following the notation of , in which one codes all the information in three variables associated with the time dependent portion of the triads $`E^1,E^2,E^3`$ and their conjugate momenta $`A_1,A_2,A_3`$. The only constraint of the theory given the homogeneity and diagonality of the metric is the Hamiltonian, which rescaled by the cosmological constant reads,

$$H=E^1E^2E^3+\mathrm{\Lambda }^1\left(E^1E^2(A_1A_2ϵA_3)+\mathrm{cyclic}\right),$$
(7)
where $`ϵ=1`$ corresponds to the Bianchi IX model and $`ϵ=0`$ to the Bianchi I model and the extra terms are obtained by cyclically permuting the indices $`1,2,3`$. It is straightforward to repeat the construction for the other class-A Bianchi models.

The observables of the zeroth order theory are given by functions $`O^0=F(E^1,E^2,E^3,E^1A_1E^2A_2,E^1A_1E^3A_3)`$ (if f one considers $`F(E^1A_1E^2A_2,E^1A_1E^3A_3)`$ these objects are exact observables to all orders for the Bianchi I case).

As an example, we can start with an observable of the zeroth order theory $`O^{(0)}=E^1`$. The condition (6) translates into the following partial differential equation

$$\frac{O^{(1)}}{A_1}E^2E^3E^1E^2A_2E^1E^3A_3+ϵE^2E^3+\mathrm{cyclic}=0.$$
(8)
This yields a first order correction, which modulo the Hamiltonian constraint can be written as,

$$O_1^{(1)}=\frac{A_1}{E^2E^3}\left(E^1E^2A_2+(E^1)^2A_1E^1E^3A_3\right),$$
(9)
up to an arbitrary additive function $`F(E^1,E^2,E^3,E^1A_1E^2A_2,E^1A_1E^3A_3)`$ that solves the homogeneous part of the equation. Expressions for higher orders can easily be generated using computer algebra, we do not list them here for reasons of space. The behavior of the approximate observables can be studied explicitly using the exact solution  of the equations of motion (in the Bianchi I case). One observes that, evaluated on the trajectories of the solution, the coefficients in the power series expansion lose their dependence on the cosmological constant as an expansion parameter. Carefully choosing the arbitrary functions $`F`$ one can define observables that are finite at the Big Bang and remain approximately constant for a significant portion of the lifetime of the (recollapsing) universe. The rate of change decreases if one adds further terms in the expansion, exhibiting convergence, although all information about $`\mathrm{\Lambda }`$ has disappeared. Therefore quite unexpectedly, the method produces solutions that are accurate irrespective of the value of the cosmological constant.

Summarizing, casting ordinary general relativity as a perturbative theory in $`\mathrm{\Lambda }^1`$ and starting with background theory in which $`\mathrm{\Lambda }\mathrm{}`$ allows to explicitly compute non-degenerate physical quantum states and observables. The technique is applicable in the theory in general and also in minisuperspace models. In the general theory, the calculations are well defined and tractable, but quite involved. This opens up a program for the quantization of the gravitational field that is well defined computationally and allows contact with the physical observables of the theory, a key ingredient in the road to finding a theory with a semiclassical regime that makes contact with ordinary low energy physics.

We wish to thank Abhay Ashtekar and Karel Kuchař for discussions. This work was supported in part by the National Science Foundation under grants PHY-9423950, INT-9811610, PHY-9407194, research funds of the Pennsylvania State University and the Eberly Family research fund at PSU. We acknowledge support of PEDECIBA.


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# The Physicist’s Guide to the Orchestra

## Introduction

Being a clarinet enthusiast, I wanted to share the richness of music (and the physicist’s way of approaching it) to engineering students within the frame of their general physics classes. The basics are reported in nearly every physics textbook and many details are provided in more specialized contributions . I could however not find a comparison of the waveforms and spectra of different instruments in any of these sources. I therefore measured and compared this data by myself.

The whole orchestra being too ambitious for a start, I restricted myself to two instrument families, namely strings and woodwinds (and the former are given a rather perfunctory treatment) and included the voice. The waveform and spectra presented in the first part illustrate the behavior of standing waves on strings and in closed and open pipes (variations 1–2), as well as highlight the basic differences between the timbre of the instruments (variations 3–5). In the second part, we will note that although the instruments are easy to identify, variations between models or performers remain hard to assess (variation 6). Furthermore, the simple models fail to explain the characteristics of some instruments (variations 6–8).

## Theme

Physicists have been studying the spectrum of musical instruments at least since the 1940s. The first available methods were based on heterodyne analysis  or on sonographs . The development of the Fast Fourier Transform (FFT) algorithm coupled with the apparition of fast and relatively cheap microprocessors has greatly facilitated the task of the musically inclined physicists. Quite sophisticated analysers have been realized  but setups based on commercial instruments work just as well for basic analysis .

The waveforms have been acquired with the setup presented in Figure 1. It consists of a condenser microphone connected directly to a 125 MHz LeCroy9400 digital oscilloscope. The spectrum was calculated by the data acquisition program LabView by FFT (this task can be directly performed on most modern oscilloscopes). On most figures, the time axis of the waveforms has been scaled and aligned for easier comparison. The spectra are given over 10 or 20 harmonics with ticks corresponding to multiples of the fundamental frequency.

The waveforms and corresponding spectra analyzed here represent the continuous part of the sound only, which is only a small part of the character of a note. The dependence on the type of attack, duration or strength was not considered, and would be a fascinating study in itself. It is also clear that a more profound analysis would require a careful consideration of instrument and microphone position, as well as a calibration of the room acoustics and of the microphone response . None of these points has been taken into account.

## Variation 1: vibrating strings

A string instrument in its simplest form is composed of a stretched string over a resonance body that transmits the vibration of the string to the air. Despite the simplicity of the vibrating system string instruments show a phenomenal diversity of timbre . This arises from the variety of excitation as the string can be plucked (guitar, harp, harpsichord), bowed (violin, viola…), or struck (piano). The resonance body plays also a great role as is attested by the timbre difference between a guitar and a banjo.

When a continuous transverse wave is generated on a stretched string of length $`l`$, a standing wave forms following the superposition of the waves after reflection at the stops. Simple considerations show that the only allowed modes of vibration correspond to wavelengths of $`\lambda =2l/n`$ where $`n1`$ , which forms a harmonic series  (inset of Figure 2). The vibration of a string will be a superposition of the different modes with varying amplitudes determined by the mode of excitation, time since the excitation etc.

One of the most simple string instrument, the string sonometer, is formed of a simple stretched string over a box-shaped resonance body. Figure 2 shows the sound produced by a plucked and bowed sonometer tuned to A2  and demonstrates the richness of the sound produced by a vibrating string as many intense upper harmonics are detected. During bowing for example, the 5th, 16th and 23rd harmonics are stronger than the fundamental.

The point and manner of excitation along the string (what physicists call the initial conditions) influences decisively the timbre: the spectra displayed in Figure 2 differ markedly, especially in the higher modes. By plucking a string, we impose an initial condition such that the shape of the string is triangular with zero and maximal displacement at the stops and at the position of the finger, respectively. The relation between position and intensity of the excited harmonics can be easily predicted (but is not easy to reproduce experimentally). This relation is not as simple for the bowed string, since the bow imparts both displacement (a fraction of mm) and velocity ($``$0.5 m/s) .

Finally, we always consider that the properties of the vibrating string are ideal. A real string has however some stiffness, which causes an increase of the frequency of the higher modes with respect to an ideal string . This can be detected on Figure 2, especially for the plucked string (presumably because of the larger displacement). As a consequence, the harmonics of a string are systematically sharp, be it for plucked , bowed  or struck strings such as in a piano .

## Variation 2: vibrating air columns

The principle of wind instruments is a bit more complicated than that of strings. The vibrating medium is the air inside a pipe that acts as a resonator where a standing wave can form. The ends of the pipe determine the boundary conditions. At a closed end, the amplitude of vibration is zero and the pressure variation is maximal (the displacement of the air and the resulting variation of pressure are in anti-phase). Conversely, the pressure will remain constant at the end of an open pipe and the standing wave shows a pressure node and a displacement antinode. This is schematically shown in the insets of Figure 4.

As a consequence, a complete series of harmonics can form in a pipe of length $`l`$ open at both ends with wavelengths equal to $`2l/n`$ with $`n1`$. If the pipe is closed at one end, the wavelength of the fundamental corresponds to four times the length of the pipe and only the odd harmonics of wavelengths equal to $`4l/(2n+1)`$ with $`n0`$ are allowed .

The vibration of the air can be excited by different means. The most simple one is to produce an edge tone by steering an airjet over an edge \[like the top of a bottle or the edge of the organ flue pipe shown in Figure 3(a)\]. The edge forms an obstacle for the jet and generates periodic vortices at the mouth of the instrument. The vortices produce in turn periodic displacements of the air molecules. When the edge forms the upper portion of a pipe, the edge tone is locked by resonance to the modes of the pipe. The pitch can then only be changed by increasing the frequency of the vortices (i.e., by blowing faster) to lock the edge tone in a higher mode (which is exactly how flautists change the register and reach higher octaves with their instruments). Such an edge excitation acts like an open end, as the vortices induce air displacement but no pressure variations.

The other means of excitation in wind instruments involve a mechanical vibration: that of the performer’s lips for brass instruments or of a reed for woodwinds. Similarly to the edge tones, the vibration of the lips or of the reed is locked to the resonances of the pipe. The simple reed of the clarinet \[see Figure 3(b)\] and of the saxophone, and the double reed of the oboe and bassoon, acts actually as a pressure regulator by admitting periodically air packets into the pipe. A reed is therefore equivalent to a closed end as it produces a pressure antinode.

We can verify the above principles with a square wooden flue organ pipe of 0.35 m length. The excitation system is reproduced in Figure 3(a) and acts as an open end, and the other end can be either closed or open. As shown on Figure 4, the fundamental of the closed pipe is found at 235 Hz, which corresponds well to a wavelength of $`\lambda =v/f=40.35=1.4`$ m with $`v=330`$ m/s. The waveform is nearly triangular, and the even harmonics are far weaker than the odd. The same pipe with its end open sounds one full octave higher (the wavelength of the fundamental is shorter by a factor of 2) and displays a complete series of harmonics.

## Variation 3: tutti

We are now ready to study the behaviour of most strings and woodwinds. Figures 5 and 6 show the waveforms and spectra of six different instruments. Table LABEL:t1 also summarizes the characteristics of the woodwinds studied here. A quick glance shows numerous disparities between the instruments, and we will try now to understand these timbre variations and their origin in more detail.

## Variation 4: the violin

The violin produces a very rich sound with at least 20 strong harmonics and complex waveforms, as was the case for the string sonometer in Figure 2. The strongest mode is not the fundamental, but the 7th harmonic in the case of Figure 6.

## Variation 5: the flute

As can be seen on Figure 5, woodwinds show simple waveforms and spectra when compared to string instruments. The flute (flauto traverso) is a textbook example of an wind instrument with open ends as the pipe is (nearly) cylindric over the whole length. The most salient feature of the flute is the limited number of harmonics ($``$7) with an intensity that decreases monotonously . The timbre is also very similar for the first two registers (not shown here).

## Variation 6 (menuetto): the clarinet

We have seen that a pipe closed at one end shows only odd harmonics: the clarinet, with its simple reed and (nearly) cylindric bore, should be a prototype of such a pipe.

At first sight, this is indeed the case. In the low register (Figure 5 for an A3 ), the odd harmonics are clearly the strongest modes. The even harmonics, although present, are strongly attenuated (at least up to the 6th harmonic). There are other marked differences with the flute. First, the sound is far richer in higher harmonics. Second, the waveform varies considerably with the pitch as displayed on Figure 7. The contents of higher harmonics strongly decreases from 20 for the A3, to 9 and 5 for the A4 and A5. The contribution of the even harmonics becomes also increasingly important. The third mode remains more intense than the second for the A4, but this is not the case anymore for A5.

The clarinet shows thus a fascinating behavior: it responds like a pipe closed at one end in the lower register but gives a sound with strong even harmonics in the higher registers. The timbre varies therefore as the pitch is increased, with a very distinctive sound for each register. This is due to several facts. First, the bore of the clarinet is not perfectly cylindric but has tapered and slightly conical sections . Second, the flared bell, the constricting mouthpiece \[Figure 3(b)\] and the toneholes (even if they are closed) perturb significantly the standing waves. Finally, for wavelengths comparable to the diameter of the toneholes, the sound wave is no longer reflected at the open tonehole but continues to propagate down the pipe. This corresponds a frequency of $``$1500 Hz in typical clarinets , and the sound will show increasing amounts of even harmonics with increasing pitch, as found on Figure 7.

Figure 7 leaves out one important feature. The clarinet does not change from the first to the second register by an octave (i.e., by doubling the frequency), but by a duodecime (tripling the frequency). This feature is due to the excitation system alone (the reed acts as a closed end), as can be easily demonstrated by replacing the mouthpiece of a flute (or of a recorder) with a clarinet mouthpiece mounted on a section of pipe such that the overall length of the instrument remains identical. The instrument sounds a full octave lower and changes registers in duodecimes, not in octaves. The reverse effect can be demonstrated by mounting a flute mouthpiece on a clarinet.

### Trio I: timbre quality

It appears from Figure 5 that it is quite easy to recognize an instrument family by its waveform or spectra. It would be tantalizing if one could also recognize one clarinet from another, for example to choose and buy a good instrument.

Figure 8 shows the spectra of my three clarinets playing the same written note , with the same mouthpiece, reed, embouchure and loudness. The upper curve correspond to my first instrument, a cheap wooden B$`\mathrm{}`$ student model. The two lower curves were obtained with professional grade B$`\mathrm{}`$ and A clarinets. I can identify each instrument by playing a few notes from the produced sound and from muscular sensations in the embouchure and respiratory apparatus.

At first glance, the spectra of the three clarinets are readily comparable. Closer inspection shows that the spectra begin to differ from the 10th harmonic on! There are actually far less variations in relative intensities between the two B$`\mathrm{}`$ instruments than between the two pro clarinets. The pro B$`\mathrm{}`$ seems to be slightly richer in harmonics than the student model. The A has no strong harmonics beyond the 11th. This leads to the conclusion that the B$`\mathrm{}`$ and A clarinets are (slightly) different instruments (many clarinetists will agree with that point). The measured differences between two B$`\mathrm{}`$ clarinets remain however quite subtle despite the huge and easily audible difference in timbre.

### Trio II: tone quality

Is it possible to tell apart a good from a bad tone? This question is of utmost importance for every musician to obtain the desired tone quality. Figure 9 shows two clarinet tones obtained on the same instrument, with the same mouthpiece and reed. The first is a good tone: one could describe it as fullbodied, agreeable to the ear. The second is a beginner’s tone: emitted with a closed throat and weak. The difference is instantly audible but difficult to quantify from Figure 9. The variations appear again in the higher harmonics: the bad tone show no harmonics beyond the 12th, which is at least five modes less than the good tone.

It is quite astonishing that the quality of the sound is determined by the presence (or absence) of high harmonics with amplitudes that are at least 40 dB (a factor 104) weaker than the fundamental! Musicians are sensitive to very subtle effects which are difficult to (a) link to a physically measurable value and (b) to quantify precisely. Conventional statistics have proven ineffective for classifying the sound quality: interestingly, effective solutions based on neural networks have been recently demonstrated .

## Variation 7: the saxophone

Can one predict the spectra of the saxophone, a single reed instrument with a truncated conical bore, by extrapolation from the previous observations? The sax is a wind instrument, which would imply a limited number of harmonics, and a spectra mainly composed of odd harmonics because of the reed. A short glance at Figures 5 and 6 shows that both predictions are wrong. The even harmonics are as strong as the odd . The sound remains very rich in harmonics even in the higher registers, far more than for the clarinet, and the timbre changes only slightly between the first and the second register. The saxophone does not behave at all like a clarinet!

The main reason is the form of the bore: in a cone, the standing waves are not plane but spherical . This has profound implications for the standing wave pattern . In short, the intensity of a wave travelling down or up the pipe is in first approximation constant along the pipe, which implies that the amplitude scales with the inverse of the distance to the cone apex. The waves interfer to form a spherical standing wave with pressure nodes separated by the usual half-wavelength spacing, but with an amplitude that varies as the inverse of the distance to the cone apex. This is true for a closed as well as an open end . A conical bore shows therefore a complete harmonic series, be it excited with a reed, the lips or an edge ! It would seem also that the conical pipe of the saxophone favors the higher harmonics as compared to the cylindric bore of the clarinet.

## Variation 8: the recorder

The predictions for the saxophone were wrong, so let’s try again with another instrument – the recorder for example. That should be easy: the bore is nearly cylindrical, it is excited by an edge and should therefore be similar to the open organ flue pipe. I expected a limited number of harmonics and a full harmonic series. Figure 5 shows that I was wrong again and this puzzled me greatly. The alto recorder indeed has a limited number of harmonics, and a similar timbre in the two registers. But it shows the spectrum of a closed pipe – the even harmonics are more suppressed than for the clarinet – and despite that it changes registers in octaves!

What is the explanation for the odd behaviour of the recorder? The player generates an airjet by blowing into a rectangular windcanal, which is then cut by the edge (see Figure 3). It appears from calculations that the position of the edge relative to the airjet influences critically the intensity of the different harmonics . When the edge cuts the side of the jet, the full harmonic series is observed. The even harmonics are however completely absent when the edge is positioned in the center of the jet, as is the case for most modern recorders (among those the one I used). This of course does not affect the modes of resonance of the instrument: the second harmonic can be excited easily by increasing the speed of the airjet, which raises the pitch by an octave. It follows also that I have been very lucky with the open organ flue pipe – which follows the expected behaviour shown on Figure 4 thanks to a favorable position of the edge with respect to the airjet !

## Variation 9: a capella

We perform frequently with a peculiar and versatile musical instrument, namely our voice. Few instruments have such varied expressive possibilities and ability to change the timbre and loudness. From the point of view of musical acoustics, the voice is a combination of a string and a wind instrument. A pipe, the vocal tract, is excited by the vibration of the vocal cords that generate a complete harmonic series as is usual for vibrating strings (see Figure 6). The timbre is however determined by the shape of the vocal tract that acts as a resonator. Depending on the position of the tongue and on the mouth opening, the position and width of the formants of the vocal tract (the broad resonances, indicated in Figure 10) can be varied and some harmonics produced by the vocal cords are favored with respect to others . Note that vocal tract and vocal cords are independent of each other, which implies that the timbre of the voice will change with the pitch for a given tract shape as the harmonics are shifted towards higher frequencies while the position of the formants remains constant.

The effect of the vocal tract shape is displayed in Figure 10 for three vowels sung at the same pitch (the formants are also indicated). The tongue is placed closed to the palate to produce the “ii”: it is nearly sinusoidal with weak upper harmonics. The first formant peaks around 200 Hz and decreases rapidly. The second and third formant around 2000 and 3000 Hz are however easily visible. The “ou” results from a single formant with a maximum around 300 Hz and a slowly decreasing tail: the waveform is more complex and richer in higher harmonics, giving a flute-like sound. The “aa” is obtained with an open tract and is far more complex. The most intense harmonic is the third because of the relatively high position ($``$800 Hz) and large width of the first formant. The second and the third formant are as intense as the first and give a significant amount of higher harmonics to the sound.

## Finale

We have seen that the physicist’s approach to musical instruments opens fascinating and complex possibilities. The classical examples (closed and open pipe, for example) are easy to reproduce, but one steps quickly into territory uncharted by the classical physics textbook, which makes the exploration all the more exciting. It remains also that instruments are easy to identify by their timbre, but that it is quite difficult to tell two different models from one another and to classify the quality of the produced sound. It may be even more difficult (not to say impossible) to examine the quality of an interpretation and to understand why well-played music touches us so deeply.

Musical acoustics is a beautiful subject to teach at every level. Music appeals to everybody and a lot of students play or have played at some stage an instrument: this makes often for lively demonstrations in front of the class. It involves both wave mechanics and fluid mechanics in quite complex ways, and a simple experimental setup can offer direct and compelling insights in the physics of sound production. I hope that this excursion in the basic physics of musical instruments will motivate some of the readers to include the subject in their curriculum and that it may provide helpful material for those who already do.

## Acknowledgments

I thank heartily the different people that either lent me their instrument or that took some time to come and play in the lab: Ariane Michellod (flute), Séverine Michellod (recorders), Stephan Fedrigo (violin – handcrafted by the performer!) and Lukas Bürgi (saxophone). I am also greatly indebted to Paul Braissant, Bernard Egger and Yvette Fazan, who maintain and expand an impressive collection of physics demonstration experiments at EPFL, and who are never put off by the sometimes strange requests of physics teachers.


text:
# Unexpected behavior in 3He films and the theory of fermions in two dimensions

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## Abstract

It is shown that the recently observed and unexplained behavior for the spin diffusion in 3He films can be understood in terms of the Fermi liquid theory, provided nonperturbative effects characteristic of the two dimensions are taken into account. It is conjectured that such a behavior suggests the onset of a regime in which spin-current is conserved.

\]

The theory of Fermi liquids in two dimensions ($`2d`$) has been at the center of various settlements involving unexpected physics that has been arising in a broad range of interesting materials. Examples range from systems in the quantum Hall regime to copper-oxide superconductors, including quantum-fluid mixtures of 3He-4He. These latter provided grounds for a recent observation of unexpected behavior, whose main cause will be argued in the following to be a consequence of the two dimensionality of a general system of fermions. In a recent Letter, Sheldon and Hallock presented experimental results for spin diffusion in a thin film of 3He on top of a superfluid 4He film. The result for the spin-diffusion coefficient ($`D`$) as a function of the 3He coverage ($`D_3`$), shows an upturn followed by a saturation and a downturn, as $`D_3`$ is increased. These features could not be reproduced by standard theoretical predictions.

Throughout their past work, Hallock and collaborators have managed to develop a notable degree of controll over the various regimes that may be set in the rich systems that constitute 3He-4He films. In the work reported in Ref., care has been taken to guarantee that the 3He film forms a $`2d`$ system of interacting fermions, with little influence from the 4He superfluid film that seats below. Hence one should be able to account for the observed features using theoretical results from a $`2d`$ Landau theory.

It is argued in this article that the theory of $`2d`$ Fermi liquids does account for the observed behavior (regarding the existence of a peak in $`D`$ as a function of $`D_3`$) provided one includes nonperturbative effects that are inherent to a Fermi liquid in two spatial dimensions. It is also pointed that the onset of regimes in the $`2d`$ Fermi liquid for which spin-current is a conserved quantity, is related to such a peaked behavior, however more experimental results are necessary to establish whether these regimes may be tuned or not.

I shall start with some remarks on existing theoretical results, and explain why the first attempts to fit the experimental data failed. I shall then devise the microscopic scenario that accounts for the physics underlying the peaked behavior observed, and calculate $`D`$, whose result is plotted in Fig.(1), compared to the experimental points. It should be stressed that the aim here is not to fully reproduce the experimental data, but yet uncover the physics that leads to the observed peak. Indeed, although the accuracy of the theoretical fitting is quite remarkable and consistent with the long withstanding experimental observations of Owers-Bradley et al. the existence of a saturation plateau is not accountable in the simple model studied here. One possible reason for such a behavior is discussed towards the end of the text.

The phenomenological result for the spin-diffusion coefficient in two dimensions was derived in Ref. and reads

$$D=\frac{\mathrm{}}{\pi m^{}}\frac{(1+F_0^a)^3C(\lambda )}{(F_0^a)^2(T_F/T)^2ln(T_F/T)},$$
(1)
where $`F_{\mathrm{}}^{a(s)}`$ stands for the usual Landau antisymmetric (symmetric) parameters, $`m^{}`$ is the effective mass, $`T_{(F)}`$ is the (Fermi) temperature, and $`C(\lambda )`$ is a function of all Landau parameters, given by Eq.(2). On fitting Eq.(1) to experimental data, one should be careful with the approximation that is taken for $`C(\lambda )`$. Based on the strictly pertubative results of Havens-Sacco and Widom, one may reasonably disregard all high order Landau parameters, which leads to an approximately constant value for $`C(\lambda )`$. Hallock and Sheldon tried to fit their experimental results using this particular approach, with $`C(\lambda )3/4`$. However, the existence of nonperturbative effects intrinsic to the $`2d`$ Fermi liquid, predicted some time ago by Engelbrecht and Randeria, suggests that such an approximation for the diffusion coefficient is not valid in general. Although single Landau parameters, which are present in the formulas of most the equilibrium properties, may not show any unusual behavior as a function of the density, the combination of various (usually infinite) Landau parameters may drive a $`2d`$ Fermi liquid into quite exotic regimes when the density is varied, as it has been pointed recently.

To carry a more specific discussion, the full expression for $`C(\lambda )`$,

$$C(\lambda )=\frac{\lambda (A_0^a)^2}{4N^2(0)z^4}\underset{n=1}{\overset{\mathrm{}}{}}\frac{(4n1)/(2n^2n)}{2n^2n1+\lambda |\mathrm{\Gamma }_{}^k(\pi )|^2},$$
(2)
is conveniently written in terms of the vertex function of perturbation theory taken in the limit of zero frequency, $`\mathrm{\Gamma }_{\sigma _1\sigma _2\sigma _3\sigma _4}^k(\theta )`$, where $`\theta `$ is the angle between the momenta of two quasiparticles in a scattering process constrained by a Fermi circle and $`\sigma _i`$ are spin labels. $`z`$ and $`N(0)`$ are respectively the Green’s function residue and the density of states, both at the Fermi surface. Here $`\lambda `$ stands for the inverse of a phase space sum,

$`\lambda ={\displaystyle \frac{8}{_{\sigma _i;\theta =0,\pi }|\mathrm{\Gamma }_{\sigma _1\sigma _2\sigma _3\sigma _4}^k(\theta )|^2\delta _{\sigma _1+\sigma _2,\sigma _3+\sigma _4}}}.`$

It is a straightforward yet not pleasant task to obtain these expressions from the results of Ref. This convenience of having closed expressions in terms of microscopic quantities is a remarkable amenability of the theory in two dimensions. The projection of $`\mathrm{\Gamma }_{\sigma _1\sigma _2\sigma _3\sigma _4}^k(\theta )`$ on a conveniently chosen basis defines the Landau scattering amplitudes $`A_{\mathrm{}}^{a(s)}`$. Using the forward scattering sum rule for these amplitudes, $`_{\mathrm{}}A_{\mathrm{}}^a+A_{\mathrm{}}^s=0`$, one readily sees that if only the leading order $`A_0^{a(s)}`$ are kept then $`C(\lambda )`$ is constant. Note, however, that varying the density may in principle drive the Landau coefficients to combine in such a way as to bring $`\mathrm{\Gamma }_{}^k(\pi )`$ down to small values, which causes $`D`$ to increase. In this particular case, if the spin-diffusion relaxation time becomes very large, either another diffusion mechanism as spin-lattice relaxation sets in or, in the ideal $`2d`$ interacting liquid, the system will enter a regime for which spin-current is conserved and a spin-viscous damping mechanism becomes associated with the transport of magnetization throughout the sample. A sharp raising of the spin-diffusion coefficient when one varies the density is then one of the signatures expected at the onset of such a regime.

In order to obtain a result that includes all the parameters, I start with the expression obtained by Engelbrecht, Randeria, and Zhang for the vertex part in the limit of zero exchange momentum,

$`\mathrm{\Gamma }_{}^\omega (\theta )={\displaystyle \frac{\pi \mathrm{}^2}{z^2m_h}}\{2\pi na_s^2`$ $`\mathrm{sin}^2(\theta /2).`$ (3)
. $`+{\displaystyle \frac{1}{\mathrm{ln}\left[2\pi na_s^2\mathrm{cos}^2(\theta /2)\right]}}\},`$ (4)
where $`m_h`$ is the 3He hydrodynamic mass, $`a_s`$ is a hard-core potential radius, and $`n`$ is the density, which is linearly related to the layer coverage $`D_3`$ in the experiment of Ref.

The expansion paramenter of the theory is $`g=(\mathrm{ln}2\pi na_s^2)^1`$, and for coverages ranging from $`00.8`$, $`g`$ falls within the interval $`00.58`$. Equation (4) follows from an all-order ladder sumation in the spirit of the long known results in three dimensions. In this sense, both terms on the right side of Eq.(4) are nonperturbative. The first term arises from a pole in the vertex function below the particle-hole continuum, and is a particularity of two dimensions, since it results from the finite value of the density of states for small momenta. This pole is associated with a two-hole bound state. Note that this pole term has an essential sigularity in $`g`$ and hence it cannot be determined to any finite order in perturbation theory. The second term contains the particle-hole continuum contributions. Although these contributions are separable in the dilute limit defined by the small values of $`g`$, the existence of a pole also phase-shifts the continuum and on determining the unusual behavior in the spin-diffusion coefficient, both terms are necessary. However, the continuum alone yields no new physics in two dimensions, it is the counterpart of the known three dimensional rendered logarithmic divergences in the repulsive channel. Indeed, a calculation using only the continuum (which will be shown elsewhere) leads to the usual behavior for $`D`$, consistent with the result of Ref. in second order perturbation, and also with Eq.(1) for a constant $`C(\lambda )`$. This discussion establishes the two dimensional nonperturbative nature of the unusual behavior in $`D`$, as derived in what follows.

In order to calculate the spin-diffusion coefficient from Eq.(1) knowledge of the zero frequency limit of the vertex part is necessary. These different limits of the vertex function are related to each other through a Bethe-Salpeter equation,

$`\mathrm{\Gamma }^k(\theta )=\mathrm{\Gamma }^\omega (\theta )N(0)z^2{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\theta ^{}}{2\pi }}\mathrm{\Gamma }^\omega (\theta \theta ^{})\mathrm{\Gamma }^k(\theta ^{}),`$ (5)
where the spin indices have been supressed for convenience. Using Eq.(4) to define the kernel plus inhomogeneous term, this equation is numerically solved and the result used to calculate $`D`$ from Eq.(1). The outcome is shown in Fig.(1) To fit the experimental data of Sheldon and Hallock at the given temperatures, a choice of $`a_s=0.75`$ Å was made. Although the width of the peak is rather sensitive to this choice, it should be noted that this value for $`a_s`$ is fairly close to the long withstanding theoretical fits to the magnetic susceptibility measuraments of Owers-Bradley et al. in 3He films on Grafoil. Such fittings were achieved by using $`a_s0.68`$ Å for a hard-core potential and $`a_s0.73`$ Å for a 6-12 potential. To account for the vertical scale, the theoretical result was multiplied by a constant fixed by making the first experimental point at $`D_3=0.1`$ for the $`30`$ mK curve coincide with the theoretical value of $`D`$ at the same coverage and temperature. The theoretical curves for all other temperatures follow with no extra adjustments. This fact, plus the consistency of the value for the hard-core radius with the experimental result for the magnetic susceptibility, which is an equilibrium response in the spin channel, are rather suggestive that the physics underlying the peaked behavior comes from the particularities of interacting fermions in a reduced geometry.

The discontinuity in the derivative of $`D`$ near $`D_3=0.8`$, and its nearly constant behavior for intermediate coverages are two features left unexplained by the above considerations. While the former is consistent with the promotion of 3He atoms to the first excited surface state of the superfluid 4He, as it is dicussed by Sheldom and Hallock in Ref. the latter needs further considerations: A full range plot of the theoretical result is shown in Fig.(2) for $`T=30`$ mK. The sharp peak shown in this figure is actually a divergence of $`D`$ occurring for $`D_30.3`$ layers. This result holds for the ideal $`2d`$ Fermi liquid and indicates that a regime for which spin-current is conserved exists for the densities near the divergence. However, the quasiparticles are coupled to the Nuclepore substrate used in the experiment. This complication is not treated here, and it causes the spin-diffusion to be several times smaller than the theoretical estimates (a mean field tortuosity factor, independent of the coverage, seems to be sufficient to correct the measured values). Hence, the spin-diffusion cannot grow indefinetly in such a non-ideal system. The fact that for intermediate values of the coverage, $`D`$ remains nearly constant as the coverage is varied seems to be consistent with the nearly constant “spin-lattice” relaxation times observed, suggesting that a saturation-like behavior might be coming from an underlying spin-relaxation mechanism with origin not in the $`2d`$ system of fermions but rather in the environment.

In summarizing, an explanation has been presented for the existence of a peak in the recent measuraments of the spin-diffusion coefficient as a function of the density in 3He films. Although the physical scenario builts from quite traditional results, it is conjectured that such a behavior might be signalizing the existence of new regimes of the $`2d`$ Fermi liquids predicted recently for which spin-current is conserved. The experimental results are shown also to support the predictions of Engelbrecht and Randeria for the existence of a collective mode in systems of interacting fermions, which exists only in two dimensions.

The author aknowledges useful discussions with K.S. Bedell and C. Kübert, the hospitality of A. Caldeira at IFGW-UNICAMP where this work has been developed, and N. Studart for quickly bringing the online version of Ref. to his attention. This work was supported in part by FAPESP, São Paulo, Brazil, Grant 98/07433-8.